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--- |
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license: cc-by-4.0 |
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pipeline_tag: text-to-image |
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tags: |
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- art |
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- ' stable-diffusion' |
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- dreambooth |
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library_name: diffusers |
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language: |
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- en |
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--- |
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<h1>This is a diffusion model for generating SD 1.5 images, trained on 11 pictures of my friend Vanya using Dreambooth.</h1> |
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<h3><i>Fine-tuned for ~14k steps using NVIDIA TESLA V100. Good similarity has been achieved despite a small dataset of 11 512 by 512 images</i></h3> |
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Usage of the token below: |
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| Token | Description | |
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|-----------------------|--------------------------------------| |
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| 👤 `VanyaKeshev` | Uses concept trained on Vanya | |
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## Examples |
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<h1>Original face sample from dataset</h1> |
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<img src="https://huggingface.co/MarkK/VanyaKeshevSD1.5/resolve/main/%D0%92%D0%B0%D0%BD%D0%B8%20%D0%9A%D1%83%D1%88%D0%B5%D0%B2%D1%8B/photo_2023-03-13_19-06-02.jpg" width="300"> |
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<h1>Result</h1> |
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<img src="https://huggingface.co/MarkK/VanyaKeshevSD1.5/resolve/main/%D0%92%D0%B0%D0%BD%D0%B8%20%D0%9A%D1%83%D1%88%D0%B5%D0%B2%D1%8B/00037-868393631.png" width="250"> |
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<img src="https://huggingface.co/MarkK/VanyaKeshevSD1.5/resolve/main/%D0%92%D0%B0%D0%BD%D0%B8%20%D0%9A%D1%83%D1%88%D0%B5%D0%B2%D1%8B/00026-1460256817.png" width="250"> |
