---
base_model: diffusers/stable-diffusion-xl-1.0-inpainting-0.1
library_name: diffusers
license: openrail++
instance_prompt: A photo of a modern-bed
widget: []
tags:
- text-to-image
- text-to-image
- diffusers-training
- diffusers
- lora
- template:sd-lora
- stable-diffusion-xl
- stable-diffusion-xl-diffusers
---

<!-- This model card has been generated automatically according to the information the training script had access to. You
should probably proofread and complete it, then remove this comment. -->


# SDXL LoRA DreamBooth - jimmy-dev/lora-trained-xl-inpainting

<Gallery />

## Model description

These are jimmy-dev/lora-trained-xl-inpainting LoRA adaption weights for diffusers/stable-diffusion-xl-1.0-inpainting-0.1.

The weights were trained  using [DreamBooth](https://dreambooth.github.io/).

LoRA for the text encoder was enabled: False.

Special VAE used for training: None.

## Trigger words

You should use A photo of a modern-bed to trigger the image generation.

## Download model

Weights for this model are available in Safetensors format.

[Download](jimmy-dev/lora-trained-xl-inpainting/tree/main) them in the Files & versions tab.



## Intended uses & limitations

#### How to use

```python
# TODO: add an example code snippet for running this diffusion pipeline
```

#### Limitations and bias

[TODO: provide examples of latent issues and potential remediations]

## Training details

[TODO: describe the data used to train the model]