4.2: IR Spectroscopy
IR Sample Preparation: A Practical Guide
Infrared spectroscopy is based on molecular vibrations caused by the oscillation of molecular dipoles. Bonds have characteristic vibrations depending on the atoms in the bond, the number of bonds and the orientation of those bonds with respect to the rest of the molecule. Thus, different molecules have specific spectra that can be collected for use in distinguishing products or identifying an unknown substance (to an extent.)
Collecting spectra through this method goes about one of three general ways. Nujol mulls and pressed pellets are typically used for collecting spectra of solids, while thin-film cells are used for solution-phase IR spectroscopy. Spectra of gases can also be obtained but will not be discussed in this guide.
Infrared Optical Materials and Handling
While it is all well and wonderful that substances can be characterized in this fashion one still has to be able to hold the substances inside of the instrument and properly prepare the samples. In an infrared spectrometer (Figure \(\PageIndex{1}\) )
the sample to be analyzed is held in front of an infrared laser beam, in order to do this, the sample must be contained in something, consequently this means that the very container the sample is in will absorb some of the infrared beam.
This is made somewhat complicated by the fact that all materials have some sort of vibration associated with them. Thus, if the sample holder has an optical window made of something that absorbs near where your sample does, the sample might not be distinguishable from the optical window of the sample holder. The range that is not blocked by a strong absorbance is known as a window (not to be confused with the optical materials of the cell.)
Windows are an important factor to consider when choosing the method to perform an analysis, as seen in Table \(\PageIndex{1}\) there are a number of different materials each with their own characteristic absorption spectra and chemical properties. Keep these factors in mind when performing analyses and precious sample will be saved. For most organic compounds NaCl works well though it is susceptible to attack from moisture. For metal coordination complexes KBr, or CsI typically work well due to their large windows. If money is not a problem then diamond or sapphire can be used for plates.
| Material | Transparent Ranges (cm -1 ) | Solubility | Notes |
|---|---|---|---|
| NaCl | 40,000 - 625 | H 2O | Easy to polish, hygroscopic |
| Silica glass | 55,000-3,000 | HF | Attacked by HF |
| Quartz | 40,000-2,500 | HF | Attacked by HF |
| Sapphire | 20,000-1,780 | - | Strong |
| Diamond | 40,000-2,500 and 1,800-200 | - | Very strong, expensive, hard, useless for pellets |
| CaF 2 | 70,000-1,110 | Acids | Attacked by acids, avoid ammonium salts |
| BaF 2 | 65,000-700 | - | Avoid ammonium salts |
| ZnSe | 10,000 - 550 | Acids | Brittle, attacked by acids |
| AgCl | 25,000-400 | - | Soft, sensitive to light. |
| KCl | 40,000-500 | H 2 O, Et 2 O, acetone | Hygroscopic, soft, easily polished, commonly used in making pellets. |
| KBr | 40,000-400 | H 2 O, EtOH | Hygroscopic, soft, easily polished, commonly used in making pellets. |
| CsBr | 10,000-250 | H 2 O, EtOH, acetone | Hygroscopic soft |
| CsI | 10,000-200 | H 2 O, EtOH, MeOH, acetone | Hygroscopic, soft. |
| Teflon | 5,000-1,200; 1,200-900 | - | Inert, disposable |
| Polyethylene | 4,000-3,000; 2,800-1,460; 1,380 - 730; 720- 30 | - | Inert, disposable |
Proper handling of these plates will ensure they have a long, useful life. Here follows a few simple pointers on how to handle plates:
- Avoid contact with solvents that the plates are soluble in.
- Keep the plates in a dessicator, the less water the better, even if the plates are insoluble to water.
- Handle with gloves, clean gloves.
- Avoid wiping the plates to prevent scratching.
That said, these simple guidelines will likely reduce most damage that can occur to a plate by simply holding it other faults such as dropping the plate from a sufficient height can result in more serious damage.
Preparation of Nujol Mulls
A common method of preparing solid samples for IR analysis is mulling. The principle here is by grinding the particles to below the wavelength of incident radiation that will be passing through there should be limited scattering. To suspend those tiny particles, an oil, often referred to as Nujol is used. IR-transparent salt plates are used to hold the sample in front of the beam in order to acquire data. To prepare a sample for IR analysis using a salt plate, first decide what segment of the frequency band should be studied, refer to Table \(\PageIndex{1}\) for the materials best suited for the sample. Figure \(\PageIndex{2}\) shows the materials needed for preparing a mull.
Preparing the mull is performed by taking a small portion of sample and adding approximately 10% of the sample volume worth of the oil and grinding this in an agate mortar and pestle as demonstrated in Figure \(\PageIndex{3}\). The resulting mull should be transparent with no visible particles.
Another method involves dissolving the solid in a solvent and allowing it to dry in the agate pestle. If using this method ensure that all of the solvent has evaporated since the solvent bands will appear in the spectrum. Some gentle heating may assist this process. This method creates very fine particles that are of a relatively consistent size. After addition of the oil further mixing (or grinding) may be necessary.
Plates should be stored in a desiccator to prevent erosion by atmospheric moisture and should appear roughly transparent. Some materials such as silicon will not, however. Gently rinse the plates with hexanes to wash any residual material off of the plates. Removing the plates from the desiccator and cleaning them should follow the preparation of the mull in order to maintain the integrity of the salt plates. Of course, if the plate is not soluble in water then it is still a good idea just to prevent the threat of mechanical trauma or a stray jet of acetone from a wash bottle.
Once the mull has been prepared, add a drop to one IR plate (Figure \(\PageIndex{4}\) ), place the second plate on top of the drop and give it a quarter turn in order to evenly coat the plate surface as seen in Figure \(\PageIndex{5}\). Place it into the spectrometer and acquire the desired data.
Always handle with gloves and preferably away from any sinks, faucets, or other sources of running or spraying water.
Spectra acquired by this method will have strong C-H absorption bands throughout several ranges 3,000 – 2,800 cm -1 and 1,500 – 1,300 cm -1 and may obscure signal.
Cleaning the plate is performed as previously mentioned with hexanes or chloroform can easily be performed by rinsing and leaving them to dry in the hood. Place the salt plates back into the desiccator as soon as reasonably possible to prevent damage. It is highly advisable to polish the plates after use, no scratches, fogging, or pits should be visible on the face of the plate. Chips, so long as they don’t cross the center of the plate are survivable but not desired. The samples of damaged salt plates in Figure \(\PageIndex{6}\) show common problems associated with use or potentially mishandling. Clouding, and to an extent, scratches can be polished out with an iron rouge. Areas where the crystal lattice is disturbed below the surface are impossible to fix and chips cannot be reattached.
FIgure \(\PageIndex{6}\) A series of plates indicating various forms of physical damage with a comparison to a good plate (Copyright: Colorado University-Boulder).
Preparation of Pellets
In an alternate method, this technique is along the same lines of the nujol mull except instead of the suspending medium being mineral oil, the suspending medium is a salt. The solid is ground into a fine powder with an agate mortar and pestle with an amount of the suspending salt. Preparing pellets with diamond for the suspending agent is somewhat illadvised considering the great hardness of the substance. Generally speaking, an amount of KBr or CsI is used for this method since they are both soft salts. Two approaches can be used to prepare pellets, one is somewhat more expensive but both usually yield decent results.
The first method is the use of a press. The salt is placed into a cylindrical holder and pressed together with a ram such as the one seen in (Figure \(\PageIndex{7}\) ). Afterwards, the pellet, in the holder, is placed into the instrument and spectra acquired.
An alternate, and cheaper method requires the use of a large hex nut with a 0.5 inch inner diameter, two bolts, and two wrenches such as the kit seen in Figure \(\PageIndex{8}\). Step-by-step instructions for loading and using the press follows:
- Screw one of the bolts into the nut about half way.
- Place the salt pellet mixture into the other opening of the nut and level by tapping the assembly on a countertop.
- Screw in the second bolt and place the assembly on its side with the bolts parallel to the countertop. Place one of the wrenches on the bolt on the right side with the handle aiming towards yourself.
- Take the second wrench and place it on the other bolt so that it attaches with an angle from the table of about 45 degrees.
- The second bolt is tightened with a body weight and left to rest for several minutes. Afterwards, the bolts are removed, and the sample placed into the instrument.
Some pellet presses also have a vacuum barb such as the one seen in (Figure \(\PageIndex{8}\). If your pellet press has one of these, consider using it as it will help remove air from the salt pellet as it is pressed. This ensures a more uniform pellet and removes absorbances in the collected spectrum due to air trapped in the pellet.
Preparation of Solution Cells
Solution cells (Figure \(\PageIndex{9}\) ) are a handy way of acquiring infrared spectra of compounds in solution and is particularly handy for monitoring reactions.
A thin-film cell consists of two salt plates with a very thin space in between them (Figure \(\PageIndex{10}\) ). Two channels allow liquid to be injected and then subsequently removed. The windows on these cells can be made from a variety of IR optical materials. One particularly useful one for water-based solutions is CaF 2 as it is not soluble in water.
Cleaning these cells can be performed by removing the solution, flushing with fresh solvent and gently removing the solvent by syringe. Do not blow air or nitrogen through the ports as this can cause mechanical deformation in the salt window if the pressure is high enough.
Deuterated Solvent Effects
One of the other aspects to solution-phase IR is that the solvent utilized in the cell has a characteristic absorption spectra. In some cases this can be alleviated by replacing the solvent with its deuterated sibling. The benefit here is that C-H bonds are now C-D bonds and have lower vibrational frequencies. Compiled in Figure \(\PageIndex{11}\) is a set of common solvents.
This effect has numerous benefits and is often applied to determining what vibrations correspond to what bond in a given molecular sample. This is often accomplished by using isotopically labeled “heavy” reagents such as ones that contain 2 H, 15 N, 18 O, or 13 C.
Basic Troubleshooting
There are numerous problems that can arise from improperly prepared samples, this section will go through some of the common problems and how to correct them. For this demonstration, spectra of ferrocene will be used. The molecular structure and a photograph of the brightly colored organometallic compound are shown in Figure \(\PageIndex{12}\) and Figure \(\PageIndex{13}\).
Figure \(\PageIndex{14}\) illustrates what a good sample of ferrocene looks like prepared in a KBr pellet. The peaks are well defined and sharp. No peak is flattened at 0% transmittance and Christiansen scattering is not evident in the baseline.
Figure \(\PageIndex{15}\) illustrates a sample with some peaks with intensities that are saturated and lose resolution making peak-picking difficult. In order to correct for this problem, scrape some of the sample off of the salt plate with a rubber spatula and reseat the opposite plate. By applying a thinner layer of sample one can improve the resolution of strongly absorbing vibrations.
Figure \(\PageIndex{16}\) illustrates a sample in which too much mineral oil was added to the mull so that the C-H bonds are far more intense than the actual sample. This can be remedied by removing the sample from the plate, grinding more sample and adding a smaller amount of the mull to the plate. Another possible way of doing this is if the sample is insoluble in hexanes, add a little to the mull and wick away the hexane-oil mixture to leave a dry solid sample. Apply a small portion of oil and replate.
Figure \(\PageIndex{17}\) illustrates the result of particles being too large and scattering light. To remedy this, remove the mull and grind further or else use the solvent deposition technique described earlier.
Characteristic IR Vibrational Modes for Hydrocarbon Compounds
| Functional group | Mode | Wavenumber range (cm -1 ) |
|---|---|---|
| CH 3 | Asymmetric stretch | 2962±10 |
| CH 3 | Symmetric stretch | 2872±10 |
| CH 3 | Asymmetric bend | 1460±10 |
| CH 3 | Symmetric bend (umbrella mode) | 1375±10 |
| CH 2 | Asymmetric stretch | 2926±10 |
| CH 2 | Symmetric stretch | 2855±10 |
| CH 2 | Scissors | 1455±10 |
| CH 2 | Rock | 720±10 |
| CH | Stretch | ~2900 (weak) |
| CH | Bend | ~1350 (weak) |
| Substitution | C-H stretch (cm -1 ) | C=C stretch (cm -1 ) | Out of plane bend (cm -1 ) |
|---|---|---|---|
| Vinyl | 3090-3075 | 1660-1630 | 900±5, 910±5 |
| Vinylidine | 3090-3075 | 1660-1630 | 890±5 |
| Cis | 3050-3000 | 1660-1630 | 690±10 |
| Trans | 3050-3000 | 1680-1665 | 965±5 |
| Tri-substituted | 3050-3000 | 1680-1665 | 815±25 |
| Tetra-substituted | - | 1680-1665 | - |
| Substitution | C-H stretch (cm -1 ) | C=C stretch (cm -1 ) | C-H wag (cm -1 ) |
|---|---|---|---|
| Mono-substituted | 3350-3250 | 2140-2100 | 700-600 |
| Di-substituted | - | 2260-2190 | - |
| Substitution | Out of plane C-H bending | Ring bend (cm -1 ) |
|---|---|---|
| Mono | 770-710 | 690±10 |
| Ortho | 810-750 | - |
| Meta | 770-735 | 690±10 |
| Para | 860-790 | - |
| Vibration | Wavenumber (cm -1 ) |
|---|---|
| CH 3 symmetric stretch | 2925±5 |
| CH 3 bend overtone | 2865±5 |
Fourier Transform Infrared Spectroscopy of Metal Ligand Complexes
The infrared (IR) range of the electromagnetic spectrum is usually divided into three regions:
- The far-infrared is always used for rotational spectroscopy, with wavenumber range 400 – 10 cm −1 and lower energy.
- The mid-infrared is suitable for a detection of the fundamental vibrations and associated rotational-vibrational structure with the frequency range approximately 4000 – 400 cm −1 .
- The near-Infrared with higher energy and wave number range 14000 – 4000 cm −1 , can excite overtone or higher harmonic vibrations.
For classical light material interaction theory, if a molecule can interact with an electromagnetic field and absorb a photon of certain frequency, the transient dipole of molecular functional group must oscillate at that frequency. Correspondingly, this transition dipole moment must be a non-zero value, however, some special vibration can be IR inactive for the stretching motion of a homonuclear diatomic molecule and vibrations do not affect the molecule’s dipole moment (e.g., N 2 ).
Mechanistic Description of the Vibrations of Polyatomic Molecules
A molecule can vibrate in many ways, and each way is called a "vibrational mode". If a molecule has N atoms, linear molecules have 3N-5 degrees of vibrational modes whereas nonlinear molecules have 3N-6 degrees of vibrational modes.
If a diatomic molecule has a harmonic vibration with the energy, \ref{1} , where n+ 1 / 2 with n = 0, 1, 2 ...). The motion of the atoms can be determined by the force equation, \ref{2} , where k is the force constant). The vibration frequency can be described by \ref{3} . In which m is actually the reduced mass (m red or μ), which is determined from the mass m 1 and m 2 of the two atoms, \ref{4} .
\[ E_{n} \ =\ -hv \label{1} \]
\[ F \ =\ -kx \label{2} \]
\[ \omega \ =\ (k/m)^{1/2} \label{3} \]
\[ m_{red} \ =\ \mu \ =\ \frac{m_{1} m_{2}}{m_{1}\ +\ m_{2} } \label{4} \]
Principle of Absorption Bands
In IR spectrum, absorption information is generally presented in the form of both wavenumber and absorption intensity or percent transmittance. The spectrum is generally showing wavenumber (cm -1 ) as the x-axis and absorption intensity or percent transmittance as the y-axis.
Transmittance, "T", is the ratio of radiant power transmitted by the sample (I) to the radiant power incident on the sample (I 0 ). Absorbance (A) is the logarithm to the base 10 of the reciprocal of the transmittance (T). The absorption intensity of molecule vibration can be determined by the Lambert-Beer Law, \label{5} . In this equation, the transmittance spectra ranges from 0 to 100%, and it can provide clear contrast between intensities of strong and weak bands. Absorbance ranges from infinity to zero. The absorption of molecules can be determined by several components. In the absorption equation, ε is called molar extinction coefficient, which is related to the molecule behavior itself, mainly the transition dipole moment, c is the concentration of the sample, and l is the sample length. Line width can be determined by the interaction with surroundings.
\[ A\ =\ log(1/T) \ =\ -log(I/I_{0} )\ =\ \varepsilon c l \label{5} \]
The Infrared Spectrometer
As shown in Figure \(\PageIndex{19}\), there are mainly four parts for fourier transform infrared spectrometer (FTIR):
- Light source. Infrared energy is emitted from a glowing black-body source as continuous radiations.
- Interferometer. It contains the interferometer, the beam splitter, the fixed mirror and the moving mirror. The beam splittertakes the incoming infrared beam and divides it into two optical beams. One beam reflects off the fixed mirror. The other beam reflects off of the moving mirror which moves a very short distance. After the divided beams are reflected from the two mirrors, they meet each other again at the beam splitter. Therefore, an interference pattern is generated by the changes in the relative position of the moving mirror to the fixed mirror. The resulting beam then passes through the sample and is eventually focused on the detector.
- Sample compartment. It is the place where the beam is transmitted through the sample. In the sample compartment, specific frequencies of energy are absorbed.
- Detector. The beam finally passes to the detector for final measurement. The two most popular detectors for a FTIR spectrometer are deuterated triglycine sulfate (pyroelectric detector) and mercury cadmium telluride (photon or quantum detector). The measured signal is sent to the computer where the Fourier transformation takes place.
A Typical Application: the detection of metal ligand complexes
Some General Absorption peaks for common types of functional groups
It is well known that all molecules chemicals have distinct absorption regions in the IR spectrum. Table \(\PageIndex{7}\) shows the absorption frequencies of common types of functional groups. For systematic evaluation, the IR spectrum is commonly divided into some sub-regions.
- In the region of 4000 - 2000 cm –1 , the appearance of absorption bands usually comes from stretching vibrations between hydrogen and other atoms. The O-H and N-H stretching frequencies range from 3700 - 3000 cm –1 . If hydrogen bond forms between O-H and other group, it generally caused peak line shape broadening and shifting to lower frequencies. The C-H stretching bands occur in the region of 3300 - 2800 cm –1 . The acetylenic C-H exhibits strong absorption at around 3300 cm –1 . Alkene and aromatic C-H stretch vibrations absorb at 3200-3000 cm –1 . Generally, asymmetric vibrational stretch frequency of alkene C-H is around 3150 cm -1 , and symmetric vibrational stretch frequency is between 3100 cm -1 and 3000 cm -1 . The saturated aliphatic C-H stretching bands range from 3000 - 2850 cm –1 , with absorption intensities that are proportional to the number of C-H bonds. Aldehydes often show two sharp C-H stretching absorption bands at 2900 - 2700 cm –1 . However, in water solution, C-H vibrational stretch is much lower than in non-polar solution. It means that the strong polarity solution can greatly reduce the transition dipole moment of C-H vibration.
- Furthermore, the stretching vibrations frequencies between hydrogen and other heteroatoms are between 2600 - 2000cm -1 , for example, S-H at 2600 - 2550 cm –1 , P-H at 2440 - 2275 cm –1 , Si-H at 2250 - 2100 cm –1 .
- The absorption bands at the 2300 - 1850 cm –1 region usually present only from triple bonds, such as C≡C at 2260 - 2100 cm –1 , C≡N at 2260 - 2000 cm –1 , diazonium salts –N≡N at approximately 2260 cm –1 , allenes C=C=C at 2000 - 1900 cm –1 . The peaks of these groups are all have strong absorption intensities. The 1950 - 1450 cm –1 region stands for double-bonded functional groups vibrational stretching.
- Most carbonyl C=O stretching bands range from 1870 - 1550 cm –1 , and the peak intensities are from mean to strong. Conjugation, ring size, hydrogen bonding, and steric and electronic effects can lead to significant shifts in absorption frequencies. Furthermore, if carbonyl links with electron-withdrawing group, such as acid chlorides and acid anhydrides, it would give rise to IR bands at 1850 - 1750 cm –1 . Ketones usually display stretching bands at 1715 cm -1 .
- None conjugated aliphatic C=C and C=N have absorption bands at 1690 - 1620 cm –1 . Besides, around 1430 and 1370cm -1 , there are two identical peaks presenting C-H bending.
- The region from 1300 - 910 cm –1 always includes the contributions from skeleton C-O and C-C vibrational stretches, giving additional molecular structural information correlated with higher frequency areas. For example, ethyl acetate not only shows its carbonyl stretch at 1750 - 1735 cm –1 , but also exhibits its identical absorption peaks at 1300 - 1000 cm –1 from the skeleton vibration of C-O and C-C stretches.
| Group | Frequency (cm -1 ) | Strength Appearance |
| C-H stretch | 2850-3400 |
Strong in nonpolar solvent
Weak in polar solvent |
| O-H stretch, N-H stretch | 3200-3700 | Broad in solvent |
|
C≡N stretch,
R-N=C=S stretch |
2050-2300 | Medium or strong |
| C≡O stretch (bound with metal) | around 2000 | Medium or strong |
| C≡C stretch | 2100-2260 | Weak |
| C=O stretch |
ca 1715 (ketone),
ca 1650 (amides) |
Strong |
| C=C stretch | 1450-1700 | Weak to strong |
| C-H bend | 1260 - 1470 | Strong |
| C-O stretch | 1040-1300 | Medium or strong |
General Introduction of Metal Ligand Complex
The metal electrons fill into the molecular orbital of ligands (CN, CO, etc.) to form complex compound. As shown in Figure \(\PageIndex{20}\), a simple molecular orbital diagram for CO can be used to explain the binding mechanism.
The CO and metal can bind with three ways:
- Donation of a pair of electrons from the C-O σ* orbital into an empty metal orbital (Figure \(\PageIndex{21}\) a).
- Donation from a metal d orbital into the C-O π* orbital to form a M-to-CO π-back bond (Figure \(\PageIndex{21}\) b).
- Under some conditions a pair of carbon π electron can donate into an empty metal d-orbital.
Some Factors to Include the Band Shifts and Strength
Herein, we mainly consider two properties: ligand stretch frequency and their absorption intensity. Take the ligand CO for example again. The frequency shift of the carbonyl peaks in the IR mainly depends on the bonding mode of the CO (terminal or bridging) and electron density on the metal. The intensity and peak numbers of the carbonyl bands depends on some factors: CO ligands numbers, geometry of the metal ligand complex and fermi resonance.
Effect on Electron Density on Metal
As shown in Table \(\PageIndex{8}\), a greater charge on the metal center result in the CO stretches vibration frequency decreasing. For example, [Ag(CO)]+show higher frequency of CO than free CO, which indicates a strengthening o
f the CO bond. σ donation removes electron density from the nonbonding HOMO of CO. From Figure, it is clear that the HOMO has a small amount of anti-bonding property, so removal of an electron actually increases (slightly) the CO bond strength. Therefore, the effect of charge and electronegativity depends on the amount of metal to CO π-back bonding and the CO IR stretching frequency.
| d x | Complex | CO stretch frequency (cm -1 ) |
| free CO | 2143 | |
| d 10 | [Ag(CO)] + | 2204 |
| d 10 | Ni(CO) 4 | 2060 |
| d 10 | [Co(CO) 4 ] - | 1890 |
| d 6 | [Mn(CO) 6 ] + | 2090 |
| d 6 | Cr(CO) 6 | 2000 |
| d 6 | [V(CO) 6 ] - | 1860 |
If the electron density on a metal center is increasing, more π-back bonding to the CO ligand(s) will also increase, as shown in Table \(\PageIndex{9}\). It means more electron density would enter into the empty carbonyl π* orbital and weaken the C-O bond. Therefore, it makes the M-CO bond strength increasing and more double-bond-like (M=C=O).
Ligation Donation Effect
Some cases, as shown in Table \(\PageIndex{9}\), different ligands would bind with same metal at the same metal-ligand complex. For example, if different electron density groups bind with Mo(CO) 3 as the same form, as shown in Figure \(\PageIndex{22}\), the CO vibrational frequencies would depend on the ligand donation effect. Compared with the PPh 3 group, CO stretching frequency which the complex binds the PF 3 group (2090, 2055 cm -1 ) is higher. It indicates that the absolute amount of electron density on that metal may have certain effect on the ability of the ligands on a metal to donate electron density to the metal center. Hence, it may be explained by the Ligand donation effect . Ligands that are trans to a carbonyl can have a large effect on the ability of the CO ligand to effectively π-backbond to the metal. For example, two trans π-backbonding ligands will partially compete for the same d-orbital electron density, weakening each other’s net M-L π-backbonding. If the trans ligand is a π-donating ligand, the free metal to CO π-backbonding can increase the M-CO bond strength (more M=C=O character). It is well known that pyridine and amines are not those strong π-donors. However, they are even worse π-backbonding ligands. So the CO is actually easy for π-back donation without any competition. Therefore, it naturally reduces the CO IR stretching frequencies in metal carbonyl complexes for the ligand donation effect.
| Metal Ligand Complex | CO Stretch Frequency (cm -1 ) |
|---|---|
| Mo(CO) 3 (PF 3 ) 3 | 2090, 2055 |
| Mo(CO) 3 [P(OMe) 3 ] 3 | 1977, 1888 |
| Mo(CO) 3 (PPh 3 ) 3 | 1934, 1835 |
| Mo(CO) 3 (NCCH 3 ) 3 | 1915, 1783 |
| Mo(CO) 3 (pyridine) 3 | 1888, 1746 |
Geometry Effects
Some cases, metal-ligand complex can form not only terminal but also bridging geometry. As shown in Figure \(\PageIndex{23}\), in the compound Fe 2 (CO) 7 (dipy), CO can act as a bridging ligand. Evidence for a bridging mode of coordination can be easily obtained through IR spectroscopy. All the metal atoms bridged by a carbonyl can donate electron density into the π* orbital of the CO and weaken the CO bond, lowering vibration frequency of CO. In this example, the CO frequency in terminal is around 2080 cm -1 , and in bridge, it shifts to around 1850 cm -1 .
Pump-probe Detection of Molecular Functional Group Vibrational Lifetime
The dynamics of molecular functional group plays an important role during a chemical process, chemical bond forming and breaking, energy transfer and other dynamics happens within picoseconds domain. It is very difficult to study such fast processes directly, for decades scientists can only learn from theoretical calculations, lacking experimental methods.
However, with the development of ultrashort pulsed laser enable experimental study of molecular functional group dynamics. With ultrafast laser technologies, people develop a series of measuring methods, among which, pump-probe technique is widely used to study the molecular functional group dynamics. Here we concentrate on how to use pump-probe experiment to measure functional group vibrational lifetime. The principle, experimental setup and data analysis will be introduced.
Principles of the Pump-probe Technique
For every function group within a molecule, such as the C≡N triple bond in phenyl selenocyanate (C 6 H 5 SeCN) or the C-D single bond in deuterated chloroform (DCCl 3 ), they have an individual infrared vibrational mode and associated energy levels. For a typical 3-level system (Figure \(\PageIndex{24}\), both the 0 to 1 and the 1 to 2 transition are near the probe pulse frequency (they don't necessarily need to have exactly the same frequency).
In a pump-probe experiment, we use the geometry as is shown in Figure \(\PageIndex{25}\). Two synchronized laser beams, one of which is called pump beam (E pu ) while the other probe beam (E pr ). There is a delay in time between each pulse. The laser pulses hit the sample, the intensity of ultrafast laser (fs or ps) is strong enough to generated 3 rd order polarization and produce 3 rd order optical response signal which is use to give dynamics information of molecular function groups. For the total response signals we have \label{6} , where µ 10 µ 21 are transition dipole moment and E 0 , E 1 , and E 2 are the energies of the three levels, and t 3 is the time delay between pump and probe beam. The delay t 3 is varied and the response signal intensity is measured. The functional group vibration life time is determined from the data.
\[ S \ =\ 4 \mu _{10} ^{4} e^{ -i(E_{1} - E_{0} ) t3/h - \Gamma t3} \label{6} \]
Typical Experimental Set-up
The optical layout of a typical pump-probe setup is schematically displayed in Figure \(\PageIndex{26}\). In the setup, the output of the oscillator (500 mW at 77 MHz repetition rate, 40 nm bandwidth centered at 800 nm) is split into two beams (1:4 power ratio). Of this, 20% of the power is to seed a femtosecond (fs) amplifier whose output is 40 fs pulses centered at 800 nm with power of ~3.4 W at 1 KHz repetition rate. The rest (80%) of the seed goes through a bandpass filter centered at 797.5nm with a width of 0.40 nm to seed a picosecond (ps) amplifier. The power of the stretched seed before entering the ps amplifier cavity is only ~3 mW. The output of the ps amplifier is 1ps pulses centered at 800 nm with a bandwidth ~0.6 nm. The power of the ps amplifier output is ~3 W. The fs amplifier is then to pump an optical parametric amplifier (OPA) which produces ~100 fs IR pulses with bandwidth of ~200 cm -1 that is tunable from 900 to 4000 cm -1 . The power of the fs IR pulses is 7~40 mW, depending on the frequencies. The ps amplifier is to pump a ps OPA which produces ~900 fs IR pulses with bandwidth of ~21 cm -1 , tunable from 900 - 4000 cm -1 . The power of the fs IR pulses is 10 ~ 40 mW, depending on frequencies.
In a typical pump-probe setup, the ps IR beam is collimated and used as the pump beam. Approximately 1% of the fs IR OPA output is used as the probe beam whose intensity is further modified by a polarizer placed before the sample. Another polarizer is placed after the sample and before the spectrograph to select different polarizations of the signal. The signal is then sent into a spectrograph to resolve frequency, and detected with a mercury cadmium telluride (MCT) dual array detector. Use of a pump pulse (femtosecond, wide band) and a probe pulse (picoseconds, narrow band), scanning the delay time and reading the data from the spectrometer, will give the lifetime of the functional group. The wide band pump and spectrometer described here is for collecting multiple group of pump-probe combination.
Data Analysis
For a typical pump-probe curve shown in Figure \(\PageIndex{27}\) life time t is defined as the corresponding time value to the half intensity as time zero.
Table \(\PageIndex{10}\) shows the pump-probe data of the C≡N triple bond in a series of aromatic cyano compounds: n -propyl cyanide (C 3 H 7 CN), ethyl thiocyanate (C 2 H 5 SCN), and ethyl selenocyanate (C 2 H 5 SeCN) for which the ν C ≡N for each compound (measured in CCl 4 solution) is 2252 cm -1 ), 2156 cm -1 , and ~2155 cm -1 , respectively.
| Delay (ps) | C 3 H 7 CN | C 2 H 5 SCN | C 2 H 5 SeCN |
|---|---|---|---|
| 0 | -0.00695 | -0.10918 | -0.06901 |
| 0.1 | -0.0074 | -0.10797 | -0.07093 |
| 0.2 | -0.00761 | -0.1071 | -0.07247 |
| 0.3 | -0.00768 | -0.10545 | -0.07346 |
| 0.4 | -0.0076 | -0.10487 | -0.07429 |
| 0.5 | -0.00778 | -0.10287 | -0.07282 |
| 0.6 | -0.00782 | -0.10286 | -0.07235 |
| 0.7 | -0.00803 | -0.10222 | -0.07089 |
| 0.8 | -0.00764 | -0.10182 | -0.07073 |
| 0.9 | -0.00776 | -0.10143 | -0.06861 |
| 1 | -0.00781 | -0.10099 | -0.06867 |
| 1.1 | -0.00745 | -0.10013 | -0.06796 |
| 1.2 | -0.00702 | -0.10066 | -0.06773 |
| 1.3 | -0.00703 | -0.0989 | -0.0676 |
| 1.4 | -0.00676 | -0.0995 | -0.06638 |
| 1.5 | -0.00681 | -0.09757 | -0.06691 |
| 1.6 | -0.00639 | -0.09758 | -0.06696 |
| 1.7 | -0.00644 | -0.09717 | -0.06583 |
| 1.8 | -0.00619 | -0.09741 | -0.06598 |
| 1.9 | -0.00613 | -0.09723 | -0.06507 |
| 2 | -0.0066 | -0.0962 | -0.06477 |
| 2.5 | -0.00574 | -0.09546 | -0.0639 |
| 3 | -0.0052 | -0.09453 | -0.06382 |
| 3.5 | -0.0482 | -0.09353 | -0.06389 |
| 4 | -0.0042 | -0.09294 | -0.06287 |
| 4.5 | -0.00387 | -0.09224 | -0.06197 |
| 5 | -0.00351 | -0.09009 | -0.06189 |
| 5.5 | -0.00362 | -0.09084 | -0.06188 |
| 6 | -0.00352 | -0.08938 | -0.06021 |
| 6.5 | -0.00269 | -0.08843 | -0.06028 |
| 7 | -0.00225 | -0.08788 | -0.05961 |
| 7.5 | -0.00231 | -0.08694 | -0.06065 |
| 8 | -0.00206 | -0.08598 | -0.05963 |
| 8.5 | -0.00233 | -0.08552 | -0.05993 |
| 9 | -0.00177 | -0.08503 | -0.05902 |
| 9.5 | -0.00186 | -0.08508 | -0.05878 |
| 10 | -0.00167 | -0.0842 | -0.0591 |
| 11 | -0.00143 | -0.08295 | -0.05734 |
A plot of intensity versus time for the data from TABLE is shown Figure \(\PageIndex{28}\). From these curves the C≡N stretch lifetimes can be determined for C 3 H 7 CN, C 2 H 5 SCN, and C 2 H 5 SeCN as ~5.5 ps, ~84 ps, and ~282 ps, respectively.
From what is shown above, the pump-probe method is used in detecting C≡N vibrational lifetimes in different chemicals. One measurement only takes several second to get all the data and the lifetime, showing that pump-probe method is a powerful way to measure functional group vibrational lifetime.
Attenuated Total Reflectace- Fourier Transform Infrared Spectroscopy
Attenuated total reflectance-Fourier transform infrared spectroscopy (ATR-FTIR) is a physical method of compositional analysis that builds upon traditional transmission FTIR spectroscopy to minimize sample preparation and optimize reproducibility. Condensed phase samples of relatively low refractive index are placed in close contact with a crystal of high refractive index and the infrared (IR) absorption spectrum of the sample can be collected. Based on total internal reflection, the absorption spectra of ATR resemble those of transmission FTIR. To learn more about transmission IR spectroscopy (FTIR) please refer to the section further up this page titled Fourier Transform Infrared Spectroscopy of Metal Ligand Complexes.
First publicly proposed in 1959 by Jacques Fahrenfort from the Royal Dutch Shell laboratories in Amsterdam, ATR IR spectroscopy was described as a technique to effectively measure weakly absorbing condensed phase materials. In Fahrenfort's first article describing the technique, published in 1961, he used a hemicylindrical ATR crystal (see Experimental Conditions) to produce single-reflection ATR (Figure \(\PageIndex{29}\) ). ATR IR spectroscopy was slow to become accepted as a method of characterization due to concerns about its quantitative effectiveness and reproducibility. The main concern being the sample and ATR crystal contact necessary to achieve decent spectral contrast. In the late 1980’s FTIR spectrometers began improving due to an increased dynamic range, signal to noise ratio, and faster computers. As a result ATR-FTIR also started gaining traction as an efficient spectroscopic technique. These days ATR accessories are often manufactured to work in conjunction with most FTIR spectrometers, as can be seen in Figure \(\PageIndex{30}\).
Total Internal Reflection
For additional information on light waves and their properties please refer to the module on Vertical Scanning Interferometry (VSI) in chapter 10.1.
When considering light propagating across an interface between two materials with different indices of refraction, the angle of refraction can be given by Snell’s law, \ref{7} , where none of the incident light will be transmitted.
\[ \varphi _{c} \ =\ \varphi _{max} \label{7} \]
The reflectance of the interface is total and whenever light is incident from a higher refractive index medium onto a lower refractive index medium, the reflection is deemed internal (as opposed to external in the opposite scenario). Total internal reflectance experiences no losses, or no transmitted light (Figure \(\PageIndex{31}\)
Supercritical internal reflection refers to angles of incidence above the critical angle of incidence allowing total internal reflectance. It is in this angular regime where only incident and reflected waves will be present. The transmitted wave is confined to the interface where its amplitude is at a maximum and will damp exponentially into the lower refractive index medium as a function of distance. This wave is referred to as the evanescent wave and it extends only a very short distance beyond the interface.
To apply total internal reflection to the experimental setup in ATR, consider n 2 to be the internal reflectance element or ATR crystal (the blue trapezoid in Figure \(\PageIndex{32}\) )
where n 2 is the material with the higher index of refraction. This should be a material that is fully transparent to the incident infrared radiation to give a real value for the refractive index. The ATR crystal must also have a high index of refraction to allow total internal reflection with many samples that have an index of refraction n 1 , where n 1 < n 2 .
We can consider the sample to be absorbing in the infrared. Electromagnetic energy will pass through the crystal/sample interface and propagate into the sample via the evanescent wave. This energy loss must be compensated with the incident IR light. Thus, total reflectance is no longer occurring and the reflection inside the crystal is attenuated. If a sample does not absorb, the reflectance at the interface shows no attenuation. Therefore if the IR light at a particular frequency does not reach the detector, the sample must have absorbed it.
The penetration depth of the evanescent wave within the sample is on the order of 1µm. The expression of the penetration depth is given in \ref{8} and is dependent upon the wavelength and angle of incident light as well as the refractive indices of the ATR crystal and sample. The effective path length is the product of the depth of penetration of the evanescent wave and the number of points that the IR light reflects at the interface between the crystal and sample. This path length is equivalent to the path length of a sample in a traditional transmission FTIR setup.
\[ d_{p} = \frac{ \lambda }{2 \pi n_{1}} (sin \omega - ( \frac{n_{1}}{n_{2}} )^{2} )^{1/2} \label{8} \]
Experimental Conditions
Refractive Indices of ATR Crystal and Sample
Typically an ATR attachment can be used with a traditional FTIR where the beam of incident IR light enters a horizontally positioned crystal with a high refractive index in the range of 1.5 to 4, as can be seen in Table \(\PageIndex{11}\) will consist of organic compounds, inorganic compounds, and polymers which have refractive indices below 2 and can readily be found on a database.
| Material | Refractive Index (RI) | Spectral Range (cm -1 ) |
|---|---|---|
| Zinc Selenide (ZnSe) | 2.4 | 20,000 - 650 |
| Germanium (Ge) | 4 | 5,500 - 870 |
| Sapphire (Al 2 O 3 ) | 1.74 | 50,000 - 2,000 |
| Diamond (C) | 2.4 |
45,000 - 2,500,
1650 - 200 |
Single and Multiple Reflection Crystals
Multiple reflection ATR was initially more popular than single reflection ATR because of the weak absorbances associated with single reflection ATR. More reflections increased the evanescent wave interaction with the sample, which was believed to increase the signal to noise ratio of the spectrum. When IR spectrometers developed better spectral contrast, single reflection ATR became more popular. The number of reflections and spectral contrast increases with the length of the crystal and decreases with the angle of incidence as well as thickness. Within multiple reflection crystals some of the light is transmitted and some is reflected as the light exits the crystal, resulting in some of the light going back through the crystal for a round trip. Therefore, light exiting the ATR crystal contains components that experienced different number of reflections at the crystal-sample interface.
Angle of Incidence
It was more common in earlier instruments to allow selection of the incident angle, sometimes offering selection between 30°, 45°, and 60°. In all cases for total internal reflection to hold, the angle of incidence must exceed the critical angle and ideally complement the angle of the crystal edge so that the light enters at a normal angle of incidence. These days 45° is the standard angle on most ATR-FTIR setups.
ATR Crystal Shape
For the most part ATR crystals will have a trapezoidal shape as shown in Figure \(\PageIndex{31}\). This shape facilitates sample preparation and handling on the crystal surface by enabling the optical setup to be placed below the crystal. However, different crystal shapes (Figure \(\PageIndex{33}\) ) may be used for particular purposes, whether it is to achieve multiple reflections or reduce the spot size. For example, a hemispherical crystal may be used in a microsampling experiment in which the beam diameter can be reduced at no expense to the light intensity. This allows appropriate measurement of a small sample without compromising the quality of the resulting spectral features.
Crystal-sample contact
Because the path length of the evanescent wave is confined to the interface between the ATR crystal and sample, the sample should make firm contact with the ATR crystal (Figure \(\PageIndex{34}\) ). The sample sits atop the crystal and intimate contact can be ensured by applying pressure above the sample. However, one must be mindful of the ATR crystal hardness. Too much pressure may distort the crystal and affect the reproducibility of the resulting spectrum.
The wavelength effect expressed in \label{7} shows an increase in penetration depth at increased wavelength. In terms of wavenumbers the relationship becomes inverse. At 4000 cm -1 penetration of the sample is 10x less than penetration at 400 cm -1 meaning the intensity of the peaks may appear higher at lower wavenumbers in the absorbance spectrum compared to the spectral features in a transmission FTIR spectrum (if an automated correction to the ATR setup is not already in place).
Selecting an ATR Crystal
ATR functions effectively on the condition that the refractive index of the crystal is of a higher refractive index than the sample. Several crystals are available for use and it is important to select an appropriate option for any given experiment (Table \(\PageIndex{11}\) ).
When selecting a material, it is important to consider reactivity, temperature, toxicity, solubility, and hardness.
The first ATR crystals in use were KRS-5, a mixture of thallium bromide and iodide, and silver halides. These materials are not listed in the table because they are not in use any longer. While cost-effective, they are not practical due to their light sensitivity, softness, and relatively low refractive indices. In addition KRS-5 is terribly toxic and dissolves on contact with many solvents, including water.
At present diamond is a favorable option for its hardness, inertness and wide spectral range, but may not be a financially viable option for some experiments. ZnSe and germanium are the most common crystal materials. ZnSe is reasonably priced, has significant mechanical strength and a long endurance. However, the surface will become etched with exposure to chemicals on either extreme of the pH scale. With a strong acid ZnSe will react to form toxic hydrogen selenide gas. ZnSe is also prone to oxidation and care must be taken to avoid the formation of an IR absorbing layer of SeO 2 . Germanium has a higher refractive index, which reduces the depth of penetration to 1 µm and may be preferable to ZnSe in applications involving intense sample absorptions or for use with samples that produce strong background absorptions. Sapphire is physically robust with a wide spectral range, but has a relatively low refractive index in terms of ATR crystals, meaning it may not be able to test as many samples as another crystal might.
Sample Versatility
Solids
The versatility of ATR is reflected in the various forms and phases that a sample can assume. Solid samples need not be compressed into a pellet, dispersed into a mull or dissolve in a solution. A ground solid sample is simply pressed to the surface of the ATR crystal. For hard samples that may present a challenge to grind into a fine solid, the total area in contact with the crystal may be compromised unless small ATR crystals with exceptional durability are used (e.g., 2 mm diamond). Loss of contact with the crystal would result in decreased signal intensity because the evanescent wave may not penetrate the sample effectively. The inherently short path length of ATR due to the short penetration depth (0.5-5 µm) enables surface-modified solid samples to be readily characterized with ATR.
Powdered samples are often tedious to prepare for analysis with transmission spectroscopy because they typically require being made into a KBr pellet to and ensuring the powdered sample is ground up sufficiently to reduce scattering. However, powdered samples require no sample preparation when taking the ATR spectra. This is advantageous in terms of time and effort, but also means the sample can easily be recovered after analysis.
Liquids
The advantage of using ATR to analyze liquid samples becomes apparent when short effective path lengths are required. The spectral reproducibility of liquid samples is certain as long as the entire length of the crystal is in contact with the liquid sample, ensuring the evanescent wave is interacting with the sample at the points of reflection, and the thickness of the liquid sample exceeds the penetration depth. A small path length may be necessary for aqueous solutions in order to reduce the absorbance of water.
Sample Preparation
ATR-FTIR has been used in fields spanning forensic analysis to pharmaceutical applications and even art preservation. Due to its ease of use and accessibility ATR can be used to determine the purity of a compound. With only a minimal amount of sample this researcher is able to collect a quick analysis of her sample and determine whether it has been adequately purified or requires further processing. As can be seen in Figure \(\PageIndex{35}\), the sample size is minute and requires no preparation. The sample is placed in close contact with the ATR crystal by turning a knob that will apply pressure to the sample (Figure \(\PageIndex{36}\) ).
ATR has an added advantage in that it inherently encloses the optical path of the IR beam. In a transmission FTIR, atmospheric compounds are constantly exposed to the IR beam and can present significant interference with the sample measurement. Of course the transmission FTIR can be purged in a dry environment, but sample measurement may become cumbersome. In an ATR measurement, however, light from the spectrometer is constantly in contact with the sample and exposure to the environment is reduced to a minimum.
Application to Inorganic Chemistry
One exciting application of ATR is in the study of classical works of art. In the study of fragments of a piece of artwork, where samples are scarce and one-of-a-kind, ATR is a suitable method of characterization because it requires only a small sample size. Determining the compounds present in art enables proper preservation and historical insight into the pieces.
In a study examining several paint samples from a various origins, a micro-ATR was employed for analysis. This study used a silicon crystal with a refractive index of 2.4 and a reduced beam size. Going beyond a simple surface analysis, this study explored the localization of various organic and inorganic compounds in the samples by performing a stratigraphic analysis. The researchers did so by embedding the samples in both KBr and a polyester resins. Two embedding techniques were compared to observe cross-sections of the samples. The mapping of the samples took approximately 1-3 hours which may seem quite laborious to some, but considering the precious nature of the sample, the wait time was acceptable to the researchers.
The optical microscope picture ( Figure \(\PageIndex{37}\) ) shows a sample of a blue painted area from the robe of a 14 th century Italian polychrome statue of a Madonna. The spectra shown in Figure \(\PageIndex{38}\) were acquired from the different layers pictured in the box marked in Figure \(\PageIndex{37}\). All spectra were collected from the cross-sectioned sample and the false-color map on each spectrum indicates the location of each of these compounds within the embedded sample. The spectra correspond to the inorganic compounds listed in Table \(\PageIndex{12}\), which also highlights characteristic vibrational bands.
| Compound | Selected Spectral Bands | Assignment |
|---|---|---|
| Cu 3 (CO 3 ) 2 (OH) 2 (Azurite) | 1493 | CO 3 2 - asymmetric stretch |
| Silicate based blue-pigments | 1035 | Si-O stretching |
| 2PbCO 3 \(\cdot \) Pb(OH) 2 (White lead) | 1399 | CO 3 2 - asymmetric stretch |
| A natural ferruginous aluminum silicate red pigment (Bole) | 3697 | OH stretching |
| CaSO 4 \(\cdot \) (Gypsum) | 1109 | SO 4 2 - asymmetric stretch |
The deep blue layer 3 corresponds to azurite and the light blue paint layer 2 to a mixture of silicate based blue pigments and white lead. Although beyond the ATR crystal’s spatial resolution limit of 20 µm, the absorption of bole was detected by the characteristic triple absorption bands of 3697, 3651, and 3619 cm -1 as seen in spectrum d of Figure \(\PageIndex{37}\). The white layer 0 was identified as gypsum.
To identify the binding material, the KBr embedded sample proved to be more effective than the polyester resin. This was due in part to the overwhelming IR absorbance of gypsum in the same spectral range (1700-1600 cm -1 ) as a characteristic stretch of the binding as well as some contaminant absorption due to the polyester embedding resin.
To spatially locate specific pigments and binding media, ATR mapping was performed on the area highlighted with a box in Figure \(\PageIndex{37}\). The false color images alongside each spectrum in Figure \(\PageIndex{38}\) indicate the relative presence of the compound corresponding to each spectrum in the boxed area. ATR mapping was achieved by taking 108 spectra across the 220x160 µm area and selecting for each identified compound by its characteristic vibrational band.


10.16: Discussion- Late Adulthood
Contributors and Attributions
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DISCUSSION : In this discussion, reflect upon and discuss BOTH of the following questions:
STEP 1 : First, write a response with at least EIGHT substantial sentences for each question, integrating concepts you learned from the reading and other materials (include links when necessary). Show that you can think critically on the topic by integrating your own thoughts, analysis, or experiences.
STEP 2 : Then return to the discussion to comment on at least TWO classmates’ posts (in at least FIVE sentences). Expand on a classmate’s comments in a value-adding, topic-related way. Promote a collaborative, supportive community, and advance the dialogue through follow-up questions. Reply posts cannot be one-liners, off-topic posts, vague statements, unsupported opinions, inadequate explanations or simply say, “I agree” or “good job.”


10.13: Molecular Term Symbols Describe Electronic States of Molecules
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Molecular term symbols specify molecular electronic energy levels. Term symbols for diatomic molecules are based on irreducible representations in linear symmetry groups, derived from spectroscopic notations. They usually consist of four parts: spin multiplicity, azimuthal angular momentum, total angular momentum and symmetry. All molecular term symbols discussed here are based on Russel-Saunders coupling .
Introduction
Molecular term symbols mark different electronic energy levels of a diatomic molecule. These symbols are similar to atomic term symbols , since both follow the Russell-Saunders coupling scheme. Molecular term symbols employ symmetry labels from group theory. The possibility of an electronic transition can be deducted from molecular term symbols following selection rules. For multi-atomic molecules, symmetry labels play most of term symbols' roles.
For homonuclear diatomics, the term symbol has the following form:
\[ ^{2S+1}\Lambda_{\Omega,(g/u)}^{(+/-)} \]
whereas Λ is the projection of the orbital angular momentum along the internuclear axis: Ω is the projection of the total angular momentum along the internuclear axis; g / u is the parity; and +/− is the reflection symmetry along an arbitrary plane containing the internuclear axis. Λ may be one of the greek letters in the sequence: Σ Π Δ Φ... when Λ = 0, 1, 2, 3..., respectively. For heteronuclear diatomics, the term symbol does not include the g / u part, for there is not inversion center in the molecule.
Determining term symbols of diatomics
Let's start with CO again. As we have seen before, the molecule has a close-shell configuration. Its ground state is a totally symmetric singlet, 1 Σ + , since the only possible values of ( S , Λ) are (0, 0). If one of the HOMO electrons on the 5σ + orbital has jumped to the LUMO, this molecule will be in an excited state as follows.
Suppose a CO molecule is in the excited state shown above. In order to know the term symbol of this state, a direct product of the labels is required for the two MO's with unpaired electrons. The multiplication is such as \(\Pi \times \Sigma^+ = \Pi\). According to Pauli's exclusion rule, these two unpaired electrons can never share the same set of quantum numbers, therefore the spin degeneracy S can reach its maximum 3. The resulting term symbols are 1 Π and 3 Π.
Now if we look at O 2 , it does not have a close-shell configuration at its ground state. There are two unpaired electrons each occupying one of the two degenerate 2 π orbitals, which can be seen in the diagram below.
The term symbol for oxygen molecule at its ground state is therefore derived such as Π x Π = Σ + + Σ - + [Δ], as the symbol in brackets does not allow the oxygen atoms to commute.
Transition between electronic states of diatomics
We'll focus on selection rules. Like atomic electronic states, different selection rules apply when differently incurred transitions occur. Usually for electric dipole field induced transitions, the selection rules are the same as for atoms.
- Δ Λ = 0, ±1 except Λ = 0 ‡ Λ ' = 0
- Δ S = 0
- Δ Ω = 0; ±1 except Ω = 0 ‡ Ω ' = 0
References
- Harris, Daniel C. (republished in 1989). Symmetry and Spectroscopy: An Introduction to Vibrational and Electronic Spectroscopy. Dover Publications. pp. 421-478. ISBN 0-486-66144-X
- D. J. Willock (2009). Molecular Symmetry . John Wiley & Sons Ltd. ISBN 0-470-85348-4


In the United States, law and government are interdependent. The US Constitution establishes the basic framework of the federal government and imposes certain limitations on the powers of government. In turn, the various branches of government are intimately involved in making, enforcing, and interpreting the law. Most law comes from Congress and the state legislatures. Courts interpret the laws and apply them to cases.
Laws are meaningless if they are not enforced. Companies have to make many decisions daily, from product development to marketing to maintaining growth. These decisions are based on financial considerations and legal requirements. If a company violates a law, it is often held accountable through litigation in courts.
Counselor’s Corner
Under the Supreme Court’s
Citizens United
decision, business entities enjoy the same right as natural persons to influence the political process through contributions. Because federal judges are appointed for life, businesses cannot directly influence actions of the judicial branch. However, they can do so indirectly by lobbying Congress on laws that it considers and lobbying the president concerning enforcement priorities. While all states have a comparable three-branch system, in some states (not in Colorado), judges obtain office through partisan elections. In such states, businesses can seek to influence the judicial branch through supporting judges whose philosophy favors business generally or a particular industry. For these reasons, in choosing whether to litigate in state or federal court, businesses should consider that federal judges may be more likely to take politically unpopular actions. ~John W., judge


3.5: Motion with Constant Acceleration (Part 1)
- Identify which equations of motion are to be used to solve for unknowns.
- Use appropriate equations of motion to solve a two-body pursuit problem.
You might guess that the greater the acceleration of, say, a car moving away from a stop sign, the greater the car’s displacement in a given time. But, we have not developed a specific equation that relates acceleration and displacement. In this section, we look at some convenient equations for kinematic relationships, starting from the definitions of displacement, velocity, and acceleration. We first investigate a single object in motion, called single-body motion. Then we investigate the motion of two objects, called two-body pursuit problems .
Notation
First, let us make some simplifications in notation. Taking the initial time to be zero, as if time is measured with a stopwatch, is a great simplification. Since elapsed time is \(\Delta\)t = t f − t 0 , taking t 0 = 0 means that \(\Delta\)t = t f , the final time on the stopwatch. When initial time is taken to be zero, we use the subscript 0 to denote initial values of position and velocity. That is, x 0 is the initial position and v 0 is the initial velocity . We put no subscripts on the final values. That is, t is the final time , x is the final position , and v is the final velocity . This gives a simpler expression for elapsed time, \(\Delta\)t = t. It also simplifies the expression for x displacement, which is now \(\Delta\)x = x − x 0 . Also, it simplifies the expression for change in velocity, which is now \(\Delta\)v = v − v 0 . To summarize, using the simplified notation, with the initial time taken to be zero,
\[\Delta t = t\]
\[\Delta x = x - x_{0}\]
\[\Delta v = v - v_{0},\]
where the subscript 0 denotes an initial value and the absence of a subscript denotes a final value in whatever motion is under consideration.
We now make the important assumption that acceleration is constant. This assumption allows us to avoid using calculus to find instantaneous acceleration. Since acceleration is constant, the average and instantaneous accelerations are equal—that is,
\[\bar{a} = a = constant \ldotp\]
Thus, we can use the symbol a for acceleration at all times. Assuming acceleration to be constant does not seriously limit the situations we can study nor does it degrade the accuracy of our treatment. For one thing, acceleration is constant in a great number of situations. Furthermore, in many other situations we can describe motion accurately by assuming a constant acceleration equal to the average acceleration for that motion. Lastly, for motion during which acceleration changes drastically, such as a car accelerating to top speed and then braking to a stop, motion can be considered in separate parts, each of which has its own constant acceleration.
Displacement and Position from Velocity
To get our first two equations, we start with the definition of average velocity:
\[\bar{v} = \frac{\Delta x}{\Delta t} \ldotp\]
Substituting the simplified notation for \(\Delta\)x and \(\Delta\)t yields
\[\bar{v} = \frac{x - x_{0}}{t} \ldotp\]
\[x = x_{0} + \bar{v} t,\label{3.10}\]
where the average velocity is
\[\bar{v} = \frac{v_{0} + v}{2} \ldotp \label{3.11}\]
The equation \(\bar{v} = \frac{v_{0} + v}{2}\) reflects the fact that when acceleration is constant, v is just the simple average of the initial and final velocities. Figure \(\PageIndex{1}\) illustrates this concept graphically. In part (a) of the figure, acceleration is constant, with velocity increasing at a constant rate. The average velocity during the 1-h interval from 40 km/h to 80 km/h is 60 km/h:
\[\bar{v} = \frac{v_{0} + v}{2} = \frac{40\; km/h + 80\; km/h}{2} = 60\; km/h \ldotp\]
In part (b), acceleration is not constant. During the 1-h interval, velocity is closer to 80 km/h than 40 km/h. Thus, the average velocity is greater than in part (a).
Solving for Final Velocity from Acceleration and Time
We can derive another useful equation by manipulating the definition of acceleration:
\[a = \frac{\Delta v}{\Delta t} \ldotp\]
Substituting the simplified notation for \(\Delta\)v and \(\Delta\)t gives us
\[a = \frac{v - v_{0}}{t}\; (constant\; a) \ldotp\]
\[v = v_{0} + at\; (constant\; a) \ldotp \label{3.12}\]
An airplane lands with an initial velocity of 70.0 m/s and then decelerates at 1.50 m/s 2 for 40.0 s. What is its final velocity?
Strategy
First, we identify the knowns: v 0 = 70 m/s, a = −1.50 m/s 2 , t = 40 s.
Second, we identify the unknown; in this case, it is final velocity v f .
Last, we determine which equation to use. To do this we figure out which kinematic equation gives the unknown in terms of the knowns. We calculate the final velocity using Equation \ref{3.12}, v = v 0 + at.
Solution
Substitute the known values and solve:
\[v = v_{0} + at = 70.0\; m/s + (-1.50\; m/s^{2})(40.0\; s) = 10.0\; m/s \ldotp\]
Figure \(\PageIndex{2}\) is a sketch that shows the acceleration and velocity vectors.
Significance
The final velocity is much less than the initial velocity, as desired when slowing down, but is still positive (see figure). With jet engines, reverse thrust can be maintained long enough to stop the plane and start moving it backward, which is indicated by a negative final velocity, but is not the case here.
In addition to being useful in problem solving, the equation v = v 0 + at gives us insight into the relationships among velocity, acceleration, and time. We can see, for example, that
- Final velocity depends on how large the acceleration is and how long it lasts
- If the acceleration is zero, then the final velocity equals the initial velocity (v = v 0 ), as expected (in other words, velocity is constant)
- If a is negative, then the final velocity is less than the initial velocity
All these observations fit our intuition. Note that it is always useful to examine basic equations in light of our intuition and experience to check that they do indeed describe nature accurately.
Solving for Final Position with Constant Acceleration
We can combine the previous equations to find a third equation that allows us to calculate the final position of an object experiencing constant acceleration. We start with
\[v = v_{0} + at \ldotp\]
Adding v 0 to each side of this equation and dividing by 2 gives
\[\frac{v_{0} + v}{2} = v_{0} + \frac{1}{2} at \ldotp\]
Since \(\frac{v_{0} + v}{2} = \bar{v}\) for constant acceleration, we have
\[\bar{v} = v_{0} + \frac{1}{2} at \ldotp\]
Now we substitute this expression for \(\bar{v}\) into the equation for displacement, x = x 0 + \(\bar{v}\)t, yielding
\[x = x_{0} + v_{0}t + \frac{1}{2} at^{2}\; (constant\; a) \ldotp \label{3.13}\]
Dragsters can achieve an average acceleration of 26.0 m/s 2 . Suppose a dragster accelerates from rest at this rate for 5.56 s Figure \(\PageIndex{3}\). How far does it travel in this time?
Strategy
First, let’s draw a sketch Figure \(\PageIndex{4}\). We are asked to find displacement, which is x if we take x 0 to be zero. (Think about x 0 as the starting line of a race. It can be anywhere, but we call it zero and measure all other positions relative to it.) We can use the equation \(x = x_{0} + v_{0}t + \frac{1}{2} at^{2}\) when we identify v 0 , a, and t from the statement of the problem.
Solution
First, we need to identify the knowns. Starting from rest means that v 0 = 0 , a is given as 26.0 m/s 2 and t is given as 5.56 s.
Second, we substitute the known values into the equation to solve for the unknown:
\[x = x_{0} + v_{0}t + \frac{1}{2} at^{2} \ldotp\]
Since the initial position and velocity are both zero, this equation simplifies to
\[x = \frac{1}{2} at^{2} \ldotp\]
Substituting the identified values of a and t gives
\[x = \frac{1}{2} (26.0\; m/s^{2})(5.56\; s)^{2} = 402\; m \ldotp\]
Significance
If we convert 402 m to miles, we find that the distance covered is very close to one-quarter of a mile, the standard distance for drag racing. So, our answer is reasonable. This is an impressive displacement to cover in only 5.56 s, but top-notch dragsters can do a quarter mile in even less time than this. If the dragster were given an initial velocity, this would add another term to the distance equation. If the same acceleration and time are used in the equation, the distance covered would be much greater.
What else can we learn by examining the equation \(x = x_{0} + v_{0}t + \frac{1}{2} at^{2}\)? We can see the following relationships:
- Displacement depends on the square of the elapsed time when acceleration is not zero. In Example 3.8, the dragster covers only one-fourth of the total distance in the first half of the elapsed time.
- If acceleration is zero, then initial velocity equals average velocity (v 0 = \(\bar{v}\)) , and \(x = x_{0} + v_{0}t + \frac{1}{2} at^{2}\) becomes x = x 0 + v 0 t.
Solving for Final Velocity from Distance and Acceleration
A fourth useful equation can be obtained from another algebraic manipulation of previous equations. If we solve v = v 0 + at for t, we get
\[t = \frac{v - v_{0}}{a} \ldotp\]
Substituting this and \(\bar{v} = \frac{v_{0} + v}{2}\) into \(x = x_{0} + \bar{v} t\), we get
\[v^{2} = v_{0}^{2} + 2a(x - x_{0})\; (constant\; a) \ldotp \label{3.14}\]
Calculate the final velocity of the dragster in Example 3.8 without using information about time.
Strategy
The equation \(v^{2} = v_{0}^{2} + 2a(x - x_{0})\) is ideally suited to this task because it relates velocities, acceleration, and displacement, and no time information is required.
Solution
First, we identify the known values. We know that v 0 = 0, since the dragster starts from rest. We also know that x − x 0 = 402 m (this was the answer in Example 3.8). The average acceleration was given by a = 26.0 m/s 2 . Second, we substitute the knowns into the equation \(v^{2} = v_{0}^{2} + 2a(x - x_{0})\) and solve for v:
\[v^{2} = 0 + 2(26.0\; m/s^{2})(402\; m) \ldotp\]
Thus,
\[v^{2} = 2.09 \times 10^{4}\; m/s^{2}\]
\[v = \sqrt{2.09 \times 10^{4}\; m^{2}/s^{2}} = 145\; m/s \ldotp\]
Significance
A velocity of 145 m/s is about 522 km/h, or about 324 mi/h, but even this breakneck speed is short of the record for the quarter mile. Also, note that a square root has two values; we took the positive value to indicate a velocity in the same direction as the acceleration.
An examination of the equation \(v^{2} = v_{0}^{2} + 2a(x - x_{0})\) can produce additional insights into the general relationships among physical quantities:
- The final velocity depends on how large the acceleration is and the distance over which it acts.
- For a fixed acceleration, a car that is going twice as fast doesn’t simply stop in twice the distance. It takes much farther to stop. (This is why we have reduced speed zones near schools.)


13.1: Establishing Protected Areas
A protected area is “a clearly defined geographical space (Figure 13.1), recognized, dedicated and managed through legal or other effective means to achieve the long-term conservation of nature with associated ecosystem services and cultural values” (Dudley, 2008). Given this broad definition, it comes as no surprise that governments, organizations, and local communities use a variety of mechanisms to establish protected areas. The most popular of these mechanisms are:
- Government action, which can occur at a national, regional, or local level.
- Community-based initiatives by local people and traditional groups.
- Land purchases and holdings by private individuals and organizations.
- Protected areas established through co-management agreements.
- Development of biological field stations or marine laboratories.
Government protected areas
Government actions are generally considered the most secure form of protection because they involve passage of laws and buy-in from multiple levels of society. Of course, legislation establishing a protected area does not guarantee that the species and ecosystems therein are adequately preserved. Small populations, especially those living in small protected areas, often require active management (Section 8.7.5) to ensure their continued survival. Another concern is that laws protecting national parks and other wildlife sanctuaries are not strictly enforced, leading to so-called paper parks—parks that appear on official government lists, but with wildlife monitoring, law enforcement, and ecosystem management lacking on the ground (Laurance et al., 2012). However, government-sanctioned protected areas do lay a solid foundation for partnerships among governments, international conservation organizations, multinational banks, research institutes, and educational organizations. Such partnerships can bring together funding, training, and scientific and management expertise to maximise the potential value of those protected areas.
Community conserved areas
In many areas, local people already protect biological communities, forests, wildlife, rivers, and coastal waters in the vicinity of their homes. Protection on these community conserved areas is enforced by village elders and councils to ensure the sustainable use of natural resources such as food supplies and drinking water. Natural areas have also been set aside by royal families and churches to provide a space for spiritual activities (see Box 2.1) and sustainable harvesting of medicinal plants (see Box 5.2). Because human activities are highly restricted in these sacred spaces, they provide an important refuge for biodiversity. Today, an increasing number of traditional communities link cultural advocacy directly to conservation through the establishment of protected areas on their lands as a safeguard against developments that would compromise their way of living. Other communities establish protected areas to attract tourists and ensure the protection of special wildlife. One such example is the Iyondji Bonobo Community Reserve in the DRC, which protects bonobos (Pan paniscus, EN), forest elephants (Loxodonta cyclotis), as well as one of the world’s most enigmatic birds, the Congo peafowl (Afropavo congensis, VU) (Dupain et al., 2013).
Traditional communities may link cultural advocacy to conservation by establishing protected areas as a safeguard against developments that would compromise their way of living.
Privately protected areas
Over the last few decades, many African countries have adopted a more Western form of land tenure under private ownership. Wealthy individuals or groups of people have taken advantage of this opportunity by acquiring large tracts of land for ecotourism purposes (de Vos et al., 2019). Because the ecotourism potential of these privately protected areas depends on how well the property is managed (Clements et al., 2016), private landowners often invest considerable effort to maintain and even increase wildlife populations on their land. Privately protected areas have unique advantages over government-protected areas. For example, they have local buy-in from landowners and their employees by design; this is often a significant stumbling block for government-protected areas. Private sites could also employ innovative funding mechanisms that allow them to fast-track land acquisition, perhaps in response to threats such as development. In some areas, privately protected areas may even employ more people, pay better wages, and contribute more to local economies that government protected areas (Sims-Castley et al., 2005). Privately protected areas can, therefore, play a significant role in overall conservation efforts (see Box 2.3), particularly in areas where threatened species (Cousins et al., 2010) and ecosystems (Gallo et al., 2009) are underrepresented in government-protected areas.
Because the ecotourism potential of private protected areas depends on how they are managed, landowners prioritize maintaining and even increasing wildlife populations on their land.
Despite the advantages of privately protected areas, we must also consider the drawbacks. Like many community conserved areas, privately protected areas are not permanently protected by the same mechanisms and oversight as government protected areas are. Ownership and management style can also change at the whim of the landowner, or perhaps the heirs. At times, management practices may be detrimental to the species and ecosystems these privately protected areas claim to protect, for example, through introduction of invasive species and harmful breeding practices (Milner et al., 2007), and by resisting regulatory controls (Cousins et al., 2010). Innovative strategies will thus be required to ensure that these areas do contribute to biodiversity protection, which include education, support, and methods that balance financial gains with conservation goals.
Co-managed protected areas
Local people who support conservation and the protection of their local natural resources are often inspired to take the lead in protecting their local biodiversity. Governments and conservation organizations can assist such initiatives by allowing local people to access specialist expertise and obtain financial assistance to develop conservation and ecotourism infrastructure. These conservation areas, characterized by partnerships between different levels of society that share decision-making responsibilities and consequences of management actions, have been termed co-managed protected areas. Tanzania, where the management of more than two million hectares of forests and woodlands have been transferred to local groups (Blomley et al., 2019), has been particularly active in this regard. One of the biggest strengths of co-management is that, with proper consultation and engagement, it avoids eco-colonialism—the unfortunate practice by some governments and conservation organizations of disregarding the rights and practices of local people during the establishment and management of new conservation areas or environmental laws and regulations.
Contractual parks offer a good model on how to avoid eco-colonialism. These protected areas are established and managed through agreements with private or communal landowners whose land forms part of a protected area (usually a national park). This not only allows a larger area to be protected, but also allows local people to benefit from biodiversity conservation through benefit sharing and job generation initiatives. Contractual parks play an important role, especially in South Africa, where it is used as a tool to meet both conservation goals and restitution of previously dispossessed land. One such example is the |Ai-|Ais/Richtersveld TFCA (Figure 13.2), which protects a huge number of succulent plant species and a variety of desert ecosystems at the border between Namibia and South Africa. Much of this national park is make up of communal lands, with the landowners—the local Nama people—having co-management and benefit-sharing agreements with the South African government (Reid et al., 2004). Incorporating activities of the landowners in sections of the park enriches tourism experiences, such as boating, hiking, and birdwatching, and contributes to preserving the Namas’ cultural identity, pastoral lifestyle, and threatened local languages (Chennels, 1999).
Field stations and marine laboratories
Biological field stations and marine laboratories are a special kind of protected area that provide a dedicated stable space for scientists, students, and even the general public to pursue research projects on all kinds of natural phenomena in an intact environment (Tydecks et al., 2016). By facilitating collaboration and long-term observation, work done at field stations in Africa has led to several fundamental scientific advances, including improved understanding of environmental responses to climate change and acid rain, as well as advances in social development through conservation activities. Today, there are biological field stations in at least 24 Sub-Saharan African countries (Tydecks et al., 2016). Among them are Namibia’s Gobabeb Research and Training Centre which focuses on desert conservation, Kenya’s Mpala Research Centre (Box 13.1) which investigates the potential for wildlife and livestock to coexist, Nigeria’s A.P. Leventis Ornithological Research Institute (see Box 15.4) which focuses on bird conservation, and Uganda’s Makerere University Biological Field Station which has a long, distinguished record of primate research.
Anchal Padukone and Dino J. Martins
Mpala Research Centre,
Nanyuki, Kenya.
In the heart of Kenya’s Laikipia district, Mpala Conservancy stretches over 200 km 2 of semi-arid savannah, acacia bushland, wooded grassland, rocky escarpments and riverine communities along the Ewaso Nyiro and Ewaso Narok rivers. The area is home to an abundance of wildlife, including all the classic savannah mammals: impala (Aepyceros melampus, LC), Grant’s gazelles (Nanger granti, LC), reticulated giraffe (Giraffa camelopardalis reticulata, EN), leopards (Panthera pardus, VU), lions (P. leo, VU), spotted hyenas (Crocuta crocuta, LC), and some of the largest savannah elephant (Loxodonta africana, VU) and African wild dog (Lycaon pictus, EN) populations in Kenya. There are also a few species typical of the northern regions of the Somali-Maasai centre of endemism, such as Grevy’s zebra (Equus grevyi, EN) and gerenuk (Litocranius walleri, NT). Mpala also functions as a working cattle ranch, with upwards of 2,000 cattle, camels, and sheep that are available for use by researchers.
This “multiple use” landscape and its neighbouring ranches provide exceptional opportunities for researchers to study interactions among humans, their domestic herds and wildlife in an area where they coexist. Since much of East Africa’s wildlife is found in similar areas outside formal protection, such research could provide essential and widely applicable knowledge for conservation efforts. They will be particularly important as conservation managers will increasingly have to balance wildlife and rangeland management needs to remain effective, in a context of human population growth and economic development. The Mpala Research Centre, established here in 1994, attracts hundreds of scientists every year who use this “living laboratory” to pursue projects varying in scope from the population biology of individual species to community-level dynamics and ecosystem functioning (Rubenstein and Rubinoff, 2014).
Having a research station in this area facilitates long-term and large-scale field experiments, including the Kenya Long-term Exclosure Experiment (KLEE). The 18 KLEE plots are designed to keep out different groups of animals: some plots only exclude megaherbivores (e.g. elephants and giraffes); others exclude all large herbivores; still others only exclude domestic cattle, among other combinations. This allows for controlled studies of the effects of different groups of herbivores on the vegetation and on each other. This research reveals that while domestic stock and wild grazers compete for forage during the dry season, the presence of zebras enhances cattle weight gain during the wet season—perhaps because zebras consume dead grass parts, improving forage quality for cattle (Riginos et al., 2012). Other studies at Mpala have also shown that wildlife and livestock can coexist and facilitate each other’s success, given the right approaches in management (Odadi et al., 2011; Ogutu et al. 2016).
Another long-time focus of research at Mpala is the threatened Grevy’s zebra. Dr. Daniel Rubenstein (Princeton University) and his research group examine the influence of environmental features on competitive behavior and reproductive patterns in plains zebra and Grevy’s zebra. In turn, they are interested in how these social processes influence zebra population size. Their findings have the potential to inform management strategies in areas where Grevy’s numbers are too low to be self-sustaining (Rubenstein, 2010).
Involving the non-scientific community, especially those living around conservancies, is crucial to the long-term success of conservation efforts. Recognising this, Mpala has hosted several citizen science initiatives. For example, The Great Grevy’s Rally was a photographic census that relied on inputs from both scientists and members of the public, who travelled to conservancies, such as Mpala, to take pictures of every Grevy’s zebra they could find. Researchers processed these images using the Image Based Ecological Information System (IBEIS, http://ibeis.org ) to differentiate individuals using their stripe patterns. This allows them to determine population size and structure, and assess whether zebra numbers are stable, increasing, or decreasing.
Also hosted at Mpala, the Kids Twiga Tally (Kahumbu et al., 2016) was a similar “sight-resight” survey of reticulated giraffes that relied on IBEIS software to distinguish between individuals and determine population structure. Its 70 young participants (Figure 13.A) came from both city schools and nearby pastoralist communities, spanning a range of socio-economic backgrounds. After spending two days taking pictures of giraffes on GPS-enabled cameras, the children returned to their schools having contributed meaningfully to conservation science.


Page 6.1: Challenges and Impacts of Energy Use
Energy for lighting, heating and cooling our buildings, manufacturing products, and powering our transportation systems comes from a variety of natural sources. The earth’s core provides geothermal energy. The gravitational pull of moon and sun create tides. The sun emits light (electromagnetic radiation), which creates wind, powers the water (hydrologic) cycle, and enables photosynthesis. Plants, algae, and cyanobacteria utilize solar energy to grow and create biomass that can be burned and used for biofuels , such as wood, biodiesel, bioethanol. Over the course of millions of years, biomass from photosynthetic organisms can create energy-rich fossil fuels through the geologic process of burial and transformation through heat and pressure.
Each of these types of energy can be defined as renewable or non-renewable. Renewable energy sources can be replenished within human lifespans. Examples include solar, wind, and biomass energy. Non-renewable energy is finite and cannot be replenished within a human timescale. Examples include nuclear energy and fossil fuels, which take millions of years to form. All energy sources have and some environmental and health cost, and the distribution of energy is not equally distributed among all nations.
Environmental and Health Challenges of Energy Use
The environmental impacts of energy use on humans and the planet can happen anywhere during the life cycle of the energy source. The impacts begin with the extraction of the resource. They continue with the processing, purification or manufacture of the source; its transportation to place of energy generation, and ends with the disposal of waste generated during use.
Extraction of fossil fuels can be used as a case study because its use has significant impacts on the environment. As we mine deeper into mountains, farther out at sea, or farther into pristine habitats, we risk damaging fragile environments, and the results of accidents or natural disasters during extraction processes can be devastating. Fossils fuels are often located far from where they are utilized so they need to be transported by pipeline, tankers, rail or trucks. These all present the potential for accidents, leakage and spills. When transported by rail or truck energy must be expended and pollutants are generated. Processing of petroleum, gas and coal generates various types of emissions and wastes, as well as utilizes water resources. Production of energy at power plants results in air, water, and, often, waste emissions. Power plants are highly regulated in the Unites States by federal and state law under the Clean Air and Clean Water Acts, while nuclear power plants are regulated by the Nuclear Regulatory Commission.
Geopolitical Challenges of Fossil Fuels
The use of fossil fuels has allowed much of the global population to reach a higher standard of living. However, this dependence on fossil fuels results in many significant impacts on society. Our modern technologies and services, such as transportation and plastics depend in many ways on fossil fuels. If supplies become limited or extremely costly, our economies are vulnerable. If countries do not have fossil fuel reserves of their own, they incur even more risk. The United States has become more and more dependent on foreign oil since 1970 when our own oil production peaked. The United States imported over half of the crude oil and refined petroleum products that we consumed during 2009. Just over half of these imports came from the Western Hemisphere (Figure \(\PageIndex{2}\)).
Figure illustrates that the United States imported over half of the crude oil and refined petroleum products that it consumed during 2009. Source: U.S. Energy Information Administration, Petroleum Supply Annual, 2009, preliminary data.The major holder of oil reserves is the Organization of Petroleum Exporting Countries, (OPEC) (Figure \(\PageIndex{3}\)). As of 2018, there were 15 member countries in OPEC: Algeria, Angola, Congo, Ecuador, Equatorial Guinea, Gabon, Iran, Iraq, Kuwait, Libya, Nigeria, Qatar, Saudi Arabia, the United Arab Emirates, and Venezuela. OPEC attempts to influence the amount of oil available to the world by assigning a production quota to each member except Iraq, for which no quota is presently set.
Figure \(\PageIndex{3}\). Proven Oil Reserves Holders Pie chart shows proven oil reserves holders. Source: C. Klein-Banai using data from BP Statistical Review of World Energy (2010)Overall compliance with these quotas is mixed since the individual countries make the actual production decisions. All of these countries have a national oil company but also allow international oil companies to operate within their borders. They can restrict the amounts of production by those oil companies. Therefore, the OPEC countries have a large influence on how much of world demand is met by OPEC and non-OPEC supply. A recent example of this is the price increases that occurred during the year 2011 after multiple popular uprisings in Arab countries, including Libya.
This pressure has lead the United States to developing policies that would reduce reliance on foreign oil such as developing additional domestic sources and obtaining it from non-Middle Eastern countries such as Canada, Mexico, Venezuela, and Nigeria. However, since fossil fuel reserves create jobs and provide dividends to investors, a lot is at stake in a nation that has oil reserves. Oil wealth may be shared with the country’s inhabitants or retained by the oil companies and dictatorships, such as in Nigeria prior to the 1990s.
Figure \(\PageIndex{4}\). Fuel Type and Carbon Emissions The two charts show the relationship between fuel type and carbon emissions for U.S. energy consumption in 2010. Source: U.S. Energy Information AdministrationContributors and Attributions
- Essentials of Environmental Science by Kamala Doršner is licensed under CC BY 4.0. Modified from the original by Matthew R. Fisher.