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10.1: R and databases
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There are many interfaces which connect R with different database management software and there is even package sqldf which allows to work with R data frames through commands from SQL language.
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10.2: R and time
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10.3: R and Bootstrapping
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All generalities like standard deviation and mean are normally taken from sample but meant to represent the whole statistical population. Therefore, it is possible that these estimations could be seriously wrong. Statistical techniques like bootstrapping were designed to minimize the risk of these errors. Bootstrap is based only on the given sample but try to estimate the whole population.
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10.4: R and shape
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Analysis of biological shape is a really useful technique. Inspired with highly influential works of D’Arcy Thompson, it takes into account not the linear measurements but the whole shape of the object: contours of teeth, bones, leaves, flower petals, and even 3D objects like skulls or beaks. Naturally, shape is not exactly measurement data, it should be analyzed with special approaches. There are methods based on the analysis of curves and methods which use landmarks and thin-plate splines (T
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10.5: R and Bayes
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10.6: R, DNA and evolution
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10.7: R and reporting
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10.8: Answers to exercises
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9.1: Introduction
The ability to reproduce in kind is a basic characteristic of all living things. In kind means that the offspring of any organism closely resembles its parent or parents. Hippopotamuses give birth to hippopotamus calves; Monterey pine trees produce seeds from which Monterey pine seedlings emerge; and adult flamingos lay eggs that hatch into flamingo chicks. In kind does not generally mean exactly the same . While many single-celled organisms and a few multicellular organisms can produce genetically identical clones of themselves through mitotic cell division, many single-celled organisms and most multicellular organisms reproduce regularly using another method.
Note that, in genetics, "parent" is often used to describe the individual organism(s) that contribute genetic material to an offspring, usually in the form of gamete cells. The concept of a genetic parent is distinct from social and legal concepts of parenthood, and may differ from those whom people consider their parents. Even within the animal kingdom, characteristics that may often be associated with sexual reproduction, such as parental care or sexual behavior, are not universal.
Sexual reproduction is the production of haploid cells and the fusion of a haploid cell from each genetic parent to form a single, unique diploid cell. In multicellular organisms, the new diploid cell will then undergo mitotic cell divisions to develop into an adult organism. A type of cell division called meiosis leads to the haploid cells that are part of the sexual reproductive cycle. Sexual reproduction, specifically meiosis and fertilization, introduces variation into offspring that may account for the evolutionary success of sexual reproduction. The vast majority of eukaryotic organisms can or must employ some form of meiosis and fertilization to reproduce.
Contributors and Attributions
Remixed and/or curated from the following works:
Fowler, S., Roush, R., & Wise, J. (2013). 7.0 Introduction . In Concepts of Biology . OpenStax (CC BY 4.0; Access for free at https://openstax.org/books/concepts-biology/pages/1-introduction ).


5.18: Assignment- Web Site Evaluation Discussion Forum
- If you have not already done so, read HOW TO EVALUATE WEB RESOURCES on the NOVA library web site.
- Choose a web site related to your field of study. For example, if you are a Psychology major, you might choose something like the American Psychological Association . If you are an Accounting major, you might choose something like The American Institute of CPAs.
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Post your responses to the following questions in the class discussion board:
- What is the URL of the web site you choose to evaluate?
- Why did you choose this web site?
- What is the purpose of the web site?
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What type of domain does it come from?
(educational, nonprofit, commercial, government, etc.) - Who published and/or wrote the pages on the web site?
- Are biases clearly stated and are affiliations clear? How do you know?
- Is the site current enough? How do you know?
- What are the author’s credentials on this subject?
- Are sources documented with footnotes or links and are there links to other resources on the topic?
- Would you recommend this web site to others? Why or why not? If you recommend it, how can others use the web site for the purpose you identified?
CC licensed content, Original
- Introduction to Computer Applications and Concepts, Website Evaluation. Authored by : Jim Shannon. Provided by : Extended Learning Institute of Northern Virginia Community College. Located at : http://www.nvcc.edu/eli/index.html . License : CC BY: Attribution


1.2.2: Affinity for Spontaneous Reaction- Chemical Potentials
For a closed system containing \(\mathrm{k}\) − chemical substances, the differential dependence of Gibbs energy on temperature, pressure and chemical composition, is given by the following equation.
\[\mathrm{dG}=-\mathrm{S} \, \mathrm{dT}+\mathrm{V} \, \mathrm{dp}+\sum_{\mathrm{j}=1}^{\mathrm{j}=\mathrm{k}}\left(\frac{\partial \mathrm{G}}{\partial \mathrm{n}_{\mathrm{j}}}\right)_{\mathrm{T}, \mathrm{p}, \mathrm{n}(\mathrm{i} \neq \mathrm{j})} \, \mathrm{dn} \mathrm{j}_{\mathrm{j}} \nonumber \]
The condition at constant \(n(i \neq j)\) indicates that the amounts of each \(i\) chemical substance except chemical substance \(j\) is constant. The Gibbs energy of a closed system is a thermodynamic potential function; equation (b).
\[\mathrm{dG}=-\mathrm{S} \, \mathrm{dT}+\mathrm{V} \, \mathrm{dp}-\mathrm{A} \, \mathrm{d} \xi \nonumber \]
Here \(\mathrm{A}\) is the affinity for spontaneous chemical reaction producing a change in extent of reaction, \(\mathrm{d} \xi\), in this case a change in composition. Further the chemical potential of chemical substance \(j\),
\[\mu_{\mathrm{j}}=\left(\frac{\partial \mathrm{G}}{\partial \mathrm{n}_{\mathrm{j}}}\right)_{\mathrm{T}, \mathrm{p}, \mathrm{n}(\mathrm{i} \neq \mathrm{j})} \nonumber \]
Comparison of equations (a) and (b) yields equation(d).
\[-A \, d \xi=\sum_{j=1}^{j=k} \mu_{j} \, d n_{j} \nonumber \]
The stoichiometry in a chemical reaction for chemical substance \(j\), \(ν_{j}\) is defined such that \(ν_{j}\) is positive for products and negative for reactants; a mnemonic is ‘P for P’.
\[v_{j}=d n_{j} / d \xi \nonumber \]
Hence the affinity for spontaneous change,
\[A=-\sum_{j=1}^{j=k} v_{j} \, \mu_{j} \nonumber \]
But at equilibrium, the affinity for spontaneous change \(\mathrm{A}\) is zero.
\[\text { Then, } \sum_{\mathrm{j}=1}^{\mathrm{j}=\mathrm{k}} \mathrm{v}_{\mathrm{j}} \, \mu_{\mathrm{j}}^{\mathrm{eq}}=0 \nonumber \]
Equation (g) in terms of its simplicity is misleading. Chemists are experts at assaying a system at equilibrium in order to determine the chemical substances present and their amounts. For example, an assay of a given system yields (for defined temperature and pressure) the amounts of un-dissociated acid \(\mathrm{CH}_{3}\mathrm{COOH}(\mathrm{aq})\), and the conjugate base \(\mathrm{CH}_{3}\mathrm{COO}^{−} (\mathrm{aq})\) and hydrogen ions at equilibrium. We write equation (g) as follows.
\[-\mu^{\mathrm{eq}}\left(\mathrm{CH}_{3} \mathrm{COOH} ; \mathrm{aq}\right)+\mu^{\mathrm{eq}}\left(\mathrm{CH}_{3} \mathrm{COO}^{-} ; \mathrm{aq}\right)+\mu^{\mathrm{eq}}\left(\mathrm{H}^{+} ; \mathrm{aq}\right)=0 \nonumber \]
Or, representing a balance of chemical potentials, (a useful approach)
\[\mu^{\mathrm{eq}}\left(\mathrm{CH}_{3} \mathrm{COOH} ; \mathrm{aq}\right)=\mu^{\mathrm{eq}}\left(\mathrm{CH}_{3} \mathrm{COO}^{-} ; \mathrm{aq}\right)+\mu^{\mathrm{eq}}\left(\mathrm{H}^{+} ; \mathrm{aq}\right) \nonumber \]
The concept of a balance of equilibrium chemical potentials at thermodynamic equilibrium is often the starting point for a description of the properties of closed systems.


2: Systems of Linear Equations
Consider the system of \(n\) linear equations and \(n\) unknowns, given by
\[\begin{aligned}a_{11}x_1+a_{12}x_2+\cdots +a_{1n}x_n&=b_1, \\ a_{21}x_1+a_{22}x_2+\cdots +a_{2n}x_n&=b_2, \\ \qquad \vdots\qquad\qquad &\vdots \\ a_{n1}x_1+a_{n2}x_2+\cdots +a_{nn}x_n&=b_n.\end{aligned} \nonumber \]
We can write this system as the matrix equation
\[\text{Ax}=\text{b},\label{eq:1} \]
with
\[\text{A}=\left(\begin{array}{cccc}a_{11}&a_{12}&\cdots &a_{1n} \\ a_{21}&a_{22}&\cdots &a_{2n} \\ \vdots&\vdots&\ddots&\vdots \\ a_{n1}&a_{n2}&\cdots&a_{nn}\end{array}\right),\quad\text{x}=\left(\begin{array}{c}x_1\\x_2\\ \vdots\\x_n\end{array}\right),\quad\text{b}=\left(\begin{array}{c}b_1\\b_2\\ \vdots\\ b_n\end{array}\right).\nonumber \]
This chapter details the standard algorithm to solve \(\eqref{eq:1}\).


15.3.3: Tsunamis
A tsunami is a fast moving, long wavelength water wave generated by a sudden displacement of water by undersea landslides, earthquakes, or volcanic activity. They are commonly called "tidal waves", but are by no means created by the tides. The sudden slippage of the ocean floor near the source of an earthquake can send a train of seismic waves across the ocean. When an overriding plate along a subduction zone suddenly breaks free it moves upward raising the sea floor and the water above. The waves move outward in ever-expanding circles, nearly imperceptible in deep ocean water. As they approach land, water recedes from the shore. This dramatic action often entices the curious to investigate. But very shortly a water wave rises rapidly from the ocean surface, crashing onto shore and rushing landward.
On December 26 th , 2004 subduction between the Indian and Eurasian (more specifically the Burma) plates off the coast off the coast of Indonesia resulted in a magnitude 9 earthquake and large tsunami that devastated South Asia. [Watch "Violent Earth" from National Geographic.] The earthquake was a result of stresses being released when the Indian plate slid beneath the Burma plate. The resulting vertical movement of the ocean floor displaced hundreds of cubic meters of water and large waves propagated outward from the focus of the quake. The massive tsunami devastated coastal regions of South Asia as it crashed ashore.
Video : "Killer Tsunamis" from National Geographic
One of the largest earthquakes to occur in recent memory occurred March 11, 2011 off the coast of Japan. The huge tsunami it created devasted the northeast coast of Japan, destroying towns and infrastructure. Learn more about this magnitude 9 earthquake and its aftermath by "Digging Deeper: The Great East Japan Earthquake of 2011", or skip and continue reading.
On March 11, 2011 the Island of Japan was rocked by a violent earthquake. The magnitude 9 Tōhoku earthquake, officially named the Great East Japan Earthquake, originated 129 km (80 miles) offshore east of Sendai , Honshu, Japan, at a depth of 30 km. It was the fourth largest earthquake in the world since 1900. Honshu lies near the junction of the North American, Eurasian, Pacific, and Phillippine Sea plates. The massive earthquake radiated from plates slipping along the subduction zone, actually moving the entire country. Areas near Ishinomaki moved southeast by 5.3 meters (17 feet), while parts of Tokyo moved 9 inches.
As the North American plate slid upward over the Pacific plate, a 16 foot bulge of water rose and swept across the Pacific at an estimated 500 mph. As the water approached land the wave increased to incredible heights. A huge tsunami washed across the northern coast of Japan within moments. Sendai airport was inundated by a wall of water 12 m(39 ft) high. An estimated tsunami height of 37.9 m (124 ft) swept across the slope of a mountain 200 m (656 ft) from the coast. Flooded fields several kilometers inland remained so after more than a week due in part to coastal subsidence. Land subsidence at Oshika Peninsula Miyagi was as much as 1.2 m (3.93 ft). Those coast locations suffering from subsidence are now more susceptible to flooding during high tides.
Video : Tōhoku earthquake tsunami propagation (Courtesy NOAA)
The effects of the tsunami were far more deadly than the earthquake itself. Nearly 15,000 people were killed and 10,000 missing. Hundreds of thousands were left homeless as whole towns were devastated. The earthquake and subsequent tsunami severely damaged the Fukushima nuclear reactor station in the towns of Okuma and Futaba , leaking radioactivity into the air and water. Large evacuations of the surrounding area and contamination of food and water supplies resulted from the leakage.
Damage from the tsunami was not confined to Japan. Tens of millions of dollars in damage to public infrastructure and private property occurred in Hawaii. Along the west coast of North America docks and harbor facilities were damaged as a tsunami surge of up to 2.4 m (8 ft) struck portions of California and Oregon. An California man attempting to photograph the rising water was swept out to sea and died.


21.7: Review Questions
A nurse is explaining a client’s lipid panel results. Which statement by the client indicates a need for further teaching?
- Answer
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“Medications can be used to increase my triglyceride levels.”
A nurse is explaining administration instructions for a client who is starting alirocumab, a PCSK9 inhibitor. Which statement will the nurse include in the instructions?
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It is self-administered as a subcutaneous injection by the client.
A client began treatment for severe hypertriglyceridemia and goes to the clinic for follow-up. Which medication does the nurse expect to find on the client’s home medication list?
- Answer
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Fenofibrate
A client calls the nurse because they experienced intolerable flushing after taking over-the-counter niacin. Which intervention might the nurse recommend?
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Take aspirin 30 minutes before taking niacin.
A nurse is evaluating a follow-up lipid panel in a client who recently began taking ezetimibe. Which of the following changes would the nurse expect to see?
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Decreased LDL-cholesterol level
A nurse educator is training new nurses on the cardiac step-down unit about medications administered on the unit. Which statement indicates understanding of the drug cholestyramine?
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It is available as an oral powder that must be mixed with fluid before taking.
A client is admitted to the cardiac unit after undergoing stent placement after an acute myocardial infarction. The provider states that the client will need high-intensity statin therapy. Which medication does the nurse expect to be ordered for the client?
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Rosuvastatin
Which of the following statements is true about pravastatin?
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It is considered a hydrophilic statin, which means it may confer a lower risk of myalgia compared with other statins.
A client with a past medical history of CAD has been taking rosuvastatin 40 mg orally daily for 1 month and has a follow-up LDL-cholesterol level of 135 mg/dL, which exceeds their goal LDL-cholesterol level of 70 mg/dL. Which would be an appropriate next step for the client?
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Add ezetimibe for combination therapy.
A client is taking atorvastatin. Which severe adverse drug reaction should the client be monitored for?
- Answer
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Rhabdomyolysis


1.1: Rhetoric is Synonymous with Empty Speech
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- Cheryl E. Ball & Drew M. Loewe ed.
- West Virginia University via Digital Publishing Institute and West Virginia University Libraries
Author: Patricia Roberts-Miller, Department of Rhetoric and Writing, University of Texas at Austin
Recently, I was at a meeting of faculty whose research and teaching interests concerned issues of environmentalism. A colleague from another department asked me what my area was. “Environmental rhetoric,” I replied. He looked slightly shocked and then commented, “Good environmentalism doesn’t have a rhetoric.” I’m in a department of rhetoric, so I teach rhetoric, read scholarly pieces on rhetoric, and attend conferences on rhetoric. However, I often forget that other faculty members’ views on rhetoric might be different than mine.
A popular view of rhetoric is that it is a straightforward model of how communication should work: A person can speak the truth simply by using words that refer to true things in the world. If she chooses not to use sentences filled with words that refer to true things in the world, then she is engaged in rhetoric. Rhetoric, in this view, is something you add on to sentences (such as metaphor) that decorates and obscures communication. If I say, “The cat is on the mat,” I am using language correctly. However, if I say, “The elegant feline languishes mournfully on the expensive carpet, waiting impatiently for what he sees as his lazy servants to open a can of salmon,” then I have added rhetoric to the first sentence, or chosen rhetoric over clear communication.
For many people, the simpler, plainer version of the sentence is not just a stylistic choice, it’s a moral one. Many people believe that the addition of more complicated words obscures the meaning of the sentence. Rhetoric, to them, is something that hides the truth. If you look at the two sentences, though, you can see that the elaborated, supposedly more rhetorical one communicates quite clearly. In fact, it communicates more effectively and precisely than “The cat is on the mat.” It might, of course, be false—there might not be such a cat; it might not be elegant; it might not be thinking much of anything; it might be quite cheerful; it might not like salmon. But the same is true of the simpler sentence—there might not be a cat; it might not be on a mat. Thus, linguistic simplicity and truthfulness aren’t necessarily connected, and linguistic complexity and truthfulness aren’t necessarily opposed.
Or, to put it another way, for a long time, philosophers of language insisted that language works by sentences having propositional content—“the cat is on the mat”—which can be expressed in various ways. Rhetoric is what we layer onto the proposition. Or, as the old saying goes, “Rhetoric is clothing on the idea.” In an Edenic world, we would all wander around naked, and we would all simply and clearly speak our thoughts; rhetoric is something we must have in this fallen world.
People who believe that rhetoric hides meaning believe that we could return to Eden by using simple, plain, and rhetoric-free language. One of several underlying assumptions is that it’s harder to lie in plain language, or that lies are more obvious when the language is less complicated. Therefore, we can trust plain language and should treat complicated language with suspicion. Oddly enough, this seemingly straightforward proposition isn’t true. In other words, this simple belief shows that an idea can be untrue and persuasive at the same time. It is also interesting that the master deceivers have generally relied on simple, yet false, claims. It’s quite likely that people believed their assertions were clear and plain and, therefore, assumed that they must be true.
The Edenic view isn’t a helpful way to think about rhetoric. It isn’t even how language works. While it’s true that the same thing can be said in different ways, there is a way of saying that thing without rhetoric. “The cat is on the mat” is still a style—the simple style—with internal rhyming and prose rhythm. It’s also structurally the rhetorical figure of chiasmus—the sentence begins and ends in an almost identical way. We can’t get away from rhetoric, but we can choose its kind.
As in all interesting arguments, it’s a question of how we’re defining terms. And rhetoric has a variety of definitions. It was first used in Platonic dialogues with very little precision. It comes from the Greek word for a person with a certain role in the Athenian Assembly (rhetor). It is believed that it was Plato who added the -ic later.
He used rhetoric in terms of speech-making as opposed to arguing in small groups.  After all, Socrates—often read as a kind of spokesman for Plato’s views— relied heavily on some fairly dodgy logical moves in the dialogues. Plato’s point seems to be that speech-making isn’t a very useful skill because making speeches to large groups (Athenian juries might have hundreds of people) is not very effective for getting to the truth. It might be effective for getting others to accept the truth one has already figured out (that seems to be the point that Socrates is making in the dialogue Phaedrus ), but, if you want to find out what’s true, argue with another individual. Do not make a speech.
Of course, Socrates makes a lot of speeches in Platonic dialogues. So, it is still murky whether or not Plato noticed the contradiction, was making a different point despite noticing the proposition, or didn’t write the dialogues to get to the truth. In fact, Plato’s overall attitude toward rhetoric is murky, even though his school, the Academy, did have rhetoric classes. They were taught by a man named Aristotle.
On the other hand, Aristotle, who was a teacher of rhetoric, neither defined rhetoric as style nor as something you add to language. He described it as a discipline and a skill that enables you to see the available means of persuasion. For Aristotle, rhetoric is about public speaking to large groups, and it is different from philosophy. So, he did share those two assumptions with Plato. But he didn’t agree with Plato about rhetoric not getting us to the truth. He thought that it could get us to the truth, but that it could also be used to deceive. It depends on the motives of the person using it.
Aristotle loved syllogisms, and seems to have believed that all reasoning could be done through them. In philosophy, to get to the truth, you try to begin with a universally valid major premise (e.g., all men are mortal). Then you have a more specific proposition related to that premise (e.g., Socrates is a man) that enables you to draw a conclusion (e.g., Socrates is mortal). But Aristotle said that this kind of reasoning doesn’t work in large assemblies for two reasons. First, during a speech, people don’t have the time to reason from universally valid major premises—if you’re arguing about whether Philip of Macedon represents a threat, it’s useless to try to find universally valid premises about tyrants or war or people from Macedon. You don’t have time. Second, the kind of things about which we make speeches—politics, ethics, military strategy, guilt or innocence, honor and dishonor—aren’t subject to certainty. There are no universally valid major premises about tyrants that will help us figure out what we need to do now and here to assess Philip. We must rely on what is probably true.
According to Aristotle, what you learn from rhetoric is how to approach political, ethical, and legal problems, how to come up with an argument when you can’t be (or, at least, shouldn’t be) certain that you’re right. You also learn how to assess other people’s arguments. Aristotle, unlike many other philosophers, doesn’t present rhetoric as an inferior discipline to philosophy (he says it’s a “counterpart”). It’s just different. It’s a pragmatic skill that helps us in decision-making.
Aristotle, being an astute observer, noticed that people argued about different things in similar ways. He came up with 28 approaches, called “lines of argument” (they’re also sometimes called “formal topoi,” which makes it seem as though they have long dresses and white ties). If I am making a speech trying to persuade people to become more active in politics, I might argue from precedent (listed as #11 of his 28 lines), or argue that the consequences of political activism are good (#13), or point out inconsistencies in the argument for political quietism (#22), and so on. Those different lines aren’t ornaments I hang on the proposition that people should be politically active; they are all different ways of thinking about the situation.
Take, for instance, Aristotle’s first line of argument: consideration of the opposite, a strategy that might structure my entire case. I might spend all my time trying to show that political activism is good because political quietism is bad. I might, however, make that just part of one speech, in which I move from how good it is to be politically active to a moving description of the tragedies associated with political quietism. Or, I might make it one paragraph, or one sentence. I might say, “Ask not what your country can do for you, but what you can do for your country.” In other words, the forms—such as consideration of the opposite—can be used to structure a clause, sentence, paragraph, speech, or (in John F. Kennedy’s case) political philosophy. Rhetoric is a way of thinking. It is not just something added to a thought derived by other means.
Does that mean that rhetoric is always good? Of course not. Rhetoric is a contingent, pragmatic, and generally (but not always) verbal way of approaching problems we face as members of communities. It is the cause as well as the consequence of thought. If we tend to think in binaries and divide everything into this versus that, then we’ll probably be drawn to the rhetorical figures that divide things into two. Continually presenting and interpreting issues in that divided way will reinforce our sense that things really are divided into two. We might then act in ways that divide things into two—we might believe that everyone is either an ally or an enemy, and thereby alienate neutral parties. Thinking and talking about everyone as ally or enemy might mean we are likely to end up in a world in which people end up treating us in that manner. Rhetoric isn’t always good, and it isn’t always bad, but it’s rarely neutral.
For instance, we might be tempted to use metaphors of disease, infection, or contamination for those groups that we don’t like. That might be a calculated decision to mislead an audience. We might not dislike the groups as vehemently as we project but we still perform for the audience to get votes, money, popularity, sales, sex, or something else. It is insincere. These types of people might make us feel unsettled and disgusted. They might even come across to us as dangerous. Thus, we call them slimy or a cancer on the body politic.  Those metaphors and that rhetoric would feel accurate, and it would convey our meaning—it is not added on; it is not ornamentation. It is what we mean. And it can hurt us as a community because it can mean that we then interpret that group’s actions through a lens of disease, threat, and danger. We can end up killing them or getting them killed because of the rhetoric we used. We can’t get away from rhetoric, but we can choose the kind of rhetoric we use.
Further Reading
For further reading on rhetoric as more than “mere rhetoric,” see especially Wayne Booth’s Rhetoric of Rhetoric , which distinguishes between “rhetrickery” and rhetoric as an inclusive method of deliberation. Eugene Garver’s Aristotle’s Rhetoric: An Art of Character is an elegant introduction to Aristotle, and Debra Hawhee and Sharon Crowley’s Ancient Rhetorics for Contemporary Students , though a textbook, explains classical and current conceptions of rhetoric usefully.
Keywords
conceptual metaphor, deliberative rhetoric, public argumentation, rhetorical topic
Author Bio
Patricia Roberts-Miller is a professor in the Department of Rhetoric and Writing at the University of Texas at Austin, where she also directs the writing center. Her scholarly and teaching interests include the history and theory of public argumentation— as she likes to put it, she’s a “scholar of train wrecks in public deliberation.” More about Trish can be found at patriciarobertsmiller.com.


18.7: (edit) Magnetic Force and Torque on a Current Loop - Motors and Meters
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Learning Objectives
By the end of this section, you will be able to:
- Evaluate the net force on a current loop in an external magnetic field
- Evaluate the net torque on a current loop in an external magnetic field
- Define the magnetic dipole moment of a current loop
Motors are the most common application of magnetic force on current-carrying wires. Motors contain loops of wire in a magnetic field. When current is passed through the loops, the magnetic field exerts torque on the loops, which rotates a shaft. Electrical energy is converted into mechanical work in the process. Once the loop’s surface area is aligned with the magnetic field, the direction of current is reversed, so there is a continual torque on the loop (Figure \(\PageIndex{1}\)). This reversal of the current is done with commutators and brushes. The commutator is set to reverse the current flow at set points to keep continual motion in the motor. A basic commutator has three contact areas to avoid and dead spots where the loop would have zero instantaneous torque at that point. The brushes press against the commutator, creating electrical contact between parts of the commutator during the spinning motion.
In a uniform magnetic field, a current-carrying loop of wire, such as a loop in a motor, experiences both forces and torques on the loop. Figure \(\PageIndex{1}\) shows a rectangular loop of wire that carries a current I and has sides of lengths a and b . The loop is in a uniform magnetic field: \(\vec{B} = B\hat{j}\). The magnetic force on a straight current-carrying wire of length l is given by \(I\vec{l} \times \vec{B}\). To find the net force on the loop, we have to apply this equation to each of the four sides. The force on side 1 is
\[\vec{F}_1 = IaB \sin(90^o - \theta) \hat{i} = IaB \cos\theta \hat{i} \nonumber \]
where the direction has been determined with the RHR-1. The current in side 3 flows in the opposite direction to that of side 1, so
\[\vec{F}_3 = -IaB \sin(90^o + \theta)\hat{i} = -IaB \cos \theta \hat{i} \nonumber \]
The currents in sides 2 and 4 are perpendicular to \(\vec{B}\) and the forces on these sides are
\[\vec{F}_2 = IbB\hat{k} \nonumber \]
\[\vec{F}_4 = -IbB\hat{k}. \nonumber \]
We can now find the net force on the loop:
\[\sum \vec{F}_{net} = \vec{F}_1 + \vec{F}_2 + \vec{F}_3 + \vec{F}_4 = 0. \nonumber \]
Although this result \((\sum F = 0)\) has been obtained for a rectangular loop, it is far more general and holds for current-carrying loops of arbitrary shapes; that is, there is no net force on a current loop in a uniform magnetic field.
To find the net torque on the current loop shown in Figure \(\PageIndex{2a}\), we first consider \(F_1\) and \(F_3\). Since they have the same line of action and are equal and opposite, the sum of their torques about any axis is zero (see Fixed-Axis Rotation ). Thus, if there is any torque on the loop, it must be furnished by \(F_2\) and \(F_4\). Let’s calculate the torques around the axis that passes through point O of Figure \(\PageIndex{2b}\) (a side view of the coil) and is perpendicular to the plane of the page. The point O is a distance x from side 2 and a distance \((a - x)\) from side 4 of the loop. The moment arms of \(F_2\) and \(F_4\) are \(x \, sin \, \theta\) and \((a - x)\space sin \, \theta\), respectively, so the net torque on the loop is
\[\sum \vec{\tau} = \vec{\tau}_1 + \vec{\tau}_2 + \vec{\tau}_3 + \vec{\tau}_4 = F_2 x \, sin \, \theta \hat{i} - F_4 (a - x) \, sin \, (\theta) \hat{i} \nonumber \] \[- IbBx \, sin \, \theta \hat{i} - IbB(a - x) sin \, \theta \hat{i}. \nonumber \]
This simplifies to \[\vec{\tau} = - IAB \, sin \, \theta \hat{i} \nonumber \] where \(A = ab\) is the area of the loop.
Notice that this torque is independent of x ; it is therefore independent of where point O is located in the plane of the current loop. Consequently, the loop experiences the same torque from the magnetic field about any axis in the plane of the loop and parallel to the x -axis.
A closed-current loop is commonly referred to as a magnetic dipole and the term IA is known as its magnetic dipole moment \(\mu\). Actually, the magnetic dipole moment is a vector that is defined as
\[\vec{\mu} = IA \hat{n} \nonumber \] where \(\hat{n}\) is a unit vector directed perpendicular to the plane of the loop (see Figure \(\PageIndex{2}\)). The direction of \(\hat{n}\) is obtained with the RHR-2—if you curl the fingers of your right hand in the direction of current flow in the loop, then your thumb points along \(\hat{n}\). If the loop contains N turns of wire, then its magnetic dipole moment is given by
\[\vec{\mu} = NIA\hat{n}. \nonumber \]
In terms of the magnetic dipole moment, the torque on a current loop due to a uniform magnetic field can be written simply as
\[\vec{\tau} = \vec{\mu} \times \vec{B}. \nonumber \]
This equation holds for a current loop in a two-dimensional plane of arbitrary shape.
Using a calculation analogous to that found in Capacitance for an electric dipole, the potential energy of a magnetic dipole is
\[U = -\vec{\mu} \cdot \vec{B}. \nonumber \]
Example \(\PageIndex{1}\): Forces and Torques on Current-Carrying Loops
A circular current loop of radius 2.0 cm carries a current of 2.0 mA. (a) What is the magnitude of its magnetic dipole moment? (b) If the dipole is oriented at 30 degrees to a uniform magnetic field of magnitude 0.50 T, what is the magnitude of the torque it experiences and what is its potential energy?
Strategy
The dipole moment is defined by the current times the area of the loop. The area of the loop can be calculated from the area of the circle. The torque on the loop and potential energy are calculated from identifying the magnetic moment, magnetic field, and angle oriented in the field.
Solution
- The magnetic moment μ is calculated by the current times the area of the loop or \(\pi r^2\). \[\mu = IA = (2.0 \times 10^{-3} A)(\pi (0.02 \, m)^2) = 2.5 \times 10^{-6} A \cdot m^2 \nonumber \]
- The torque and potential energy are calculated by identifying the magnetic moment, magnetic field, and the angle between these two vectors. The calculations of these quantities are: \[\tau = \vec{\mu} \times \vec{B} = \mu B \, sin \, \theta = (2.5 \times 10^{-6} A \cdot m^2)(0.50 T) sin(30^o) = 6.3 \times 10^{-7}N \cdot m \nonumber \] \[U = -\vec{\mu} \cdot \vec{B} = - \mu B cos \theta = - (2.5 \times 10^{-6} A \cdot m^2)(0.50 T) cos (30^o) = -1.1 \times 10^{-6}J. \nonumber \]
Significance
The concept of magnetic moment at the atomic level is discussed in the next chapter. The concept of aligning the magnetic moment with the magnetic field is the functionality of devices like magnetic motors, whereby switching the external magnetic field results in a constant spinning of the loop as it tries to align with the field to minimize its potential energy.
Exercise \(\PageIndex{1}\)
In what orientation would a magnetic dipole have to be to produce (a) a maximum torque in a magnetic field? (b) A maximum energy of the dipole?
a. aligned or anti-aligned; b. perpendicular


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32: Impressionism and Extended Tonality
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32.1: Impressionism
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Impressionism is associated with Claude Debussy and Maurice Ravel in France, Ottorino Respighi in Italy, Charles Tomlinson Griffes in America, and Frederick Delius in England. We will focus on just three techniques found in the music of Debussy and Ravel: (1) the use of modes, (2) the use of upper extensions above the 7th in chord construction in tertian harmonies (9ths, 11ths, and 13ths, which we saw in the chapter on jazz), and (3) parallelism, also known as “planing.”
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32.2: Pandiatonicism
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Pandiatonicism refers to the use of all diatonic notes without the need for scale degrees or harmonies to progress or function tonally (V doesn’t need to progress to I, 7^ doesn’t need to resolve to 8^ , etc.). One often hears it as a wash of notes from the major scale, or as chords made of non-traditional combinations of notes from a major scale, often with at least one interval of a 2nd in a chord voicing.
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32.3: Quartal, Quintal, and Secundal Harmony
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Quartal harmony refers to chords stacked entirely (or mostly) in fourths. Quintal harmony refers to chords stacked entirely (or mostly) in fifths. Secundal harmony refers to chords stacked entirely (or mostly) in seconds. Contrast these concepts with the tertian (stacked in thirds) harmony we’ve encountered in the majority of this text. Alternative methods of stacking chords became more common after 1900 as composers sought ways to innovate and break with the past.
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32.4: Polychords
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A polychord typically consists of two triads sounding simultaneously. A polychord could also consist of two seventh chords, or a seventh chord and triad. Additionally, a polychord could conceivably consist of more than two triads or seventh chords, since the prefix “poly” means “many.”
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32.5: Practice Exercises
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5.8: Reading- IP Address
Introduction
An Internet Protocol address ( IP address ) is a numerical label assigned to each device (e.g., computer, printer) participating in a computer network that uses the Internet Protocol for communication. An IP address serves two principal functions: host or network interface identification and location addressing. Its role has been characterized as follows: “A name indicates what we seek. An address indicates where it is. A route indicates how to get there.”
The designers of the Internet Protocol defined an IP address as a 32-bit number and this system, known as Internet Protocol Version 4 (IPv4), is still in use today. However, because of the growth of the Internet and the predicted depletion of available addresses, a new version of IP (IPv6), using 128 bits for the address, was developed in 1995. IPv6 was standardized as RFC 2460 in 1998, and its deployment has been ongoing since the mid-2000s.
IP addresses are usually written and displayed in human-readable notations, such as <IP_ADDRESS> (IPv4), and 2001:db8:0:1234:0:567:8:1 (IPv6).
The Internet Assigned Numbers Authority (IANA) manages the IP address space allocations globally and delegates five regional Internet registries (RIRs) to allocate IP address blocks to local Internet registries (Internet service providers) and other entities.
Contents
IP versions
Two versions of the Internet Protocol (IP) are in use: IP Version 4 and IP Version 6. Each version defines an IP address differently. Because of its prevalence, the generic term IP address typically still refers to the addresses defined by IPv4. The gap in version sequence between IPv4 and IPv6 resulted from the assignment of number 5 to the experimental Internet Stream Protocol in 1979, which however was never referred to as IPv5.
IPv4 addresses
In IPv4 an address consists of 32 bits which limits the address space to<PHONE_NUMBER> (2 32 ) possible unique addresses. IPv4 reserves some addresses for special purposes such as private networks (~18 million addresses) or multicast addresses (~270 million addresses).
IPv4 addresses are canonically represented in dot-decimal notation, which consists of four decimal numbers, each ranging from 0 to 255, separated by dots, e.g., <IP_ADDRESS>. Each part represents a group of 8 bits (octet) of the address. In some cases of technical writing, IPv4 addresses may be presented in various hexadecimal, octal, or binary representations.
Subnetting
In the early stages of development of the Internet Protocol, network administrators interpreted an IP address in two parts: network number portion and host number portion. The highest order octet (most significant eight bits) in an address was designated as the network number and the remaining bits were called the rest field or host identifier and were used for host numbering within a network.
This early method soon proved inadequate as additional networks developed that were independent of the existing networks already designated by a network number. In 1981, the Internet addressing specification was revised with the introduction of classful network architecture.
Classful network design allowed for a larger number of individual network assignments and fine-grained subnetwork design. The first three bits of the most significant octet of an IP address were defined as the class of the address. Three classes ( A , B , and C ) were defined for universal unicast addressing. Depending on the class derived, the network identification was based on octet boundary segments of the entire address. Each class used successively additional octets in the network identifier, thus reducing the possible number of hosts in the higher order classes ( B and C ). The following table gives an overview of this now obsolete system.
| Class |
Leading
bits |
Size of
network
number bit field |
Size of
rest
bit field |
Number
of networks |
Addresses
per network |
Start address | End address |
|---|---|---|---|---|---|---|---|
| A | 0 | 8 | 24 | 128 (2 7 ) | 16,777,216 (2 24 ) | <IP_ADDRESS> | <IP_ADDRESS> |
| B | 10 | 16 | 16 | 16,384 (2 14 ) | 65,536 (2 16 ) | <IP_ADDRESS> | <IP_ADDRESS> |
| C | 110 | 24 | 8 | 2,097,152 (2 21 ) | 256 (2 8 ) | <IP_ADDRESS> | <IP_ADDRESS> |
Classful network design served its purpose in the startup stage of the Internet, but it lacked scalability in the face of the rapid expansion of the network in the 1990s. The class system of the address space was replaced with Classless Inter-Domain Routing (CIDR) in 1993. CIDR is based on variable-length subnet masking (VLSM) to allow allocation and routing based on arbitrary-length prefixes.
Today, remnants of classful network concepts function only in a limited scope as the default configuration parameters of some network software and hardware components (e.g. netmask), and in the technical jargon used in network administrators’ discussions.
Private addresses
Early network design, when global end-to-end connectivity was envisioned for communications with all Internet hosts, intended that IP addresses be uniquely assigned to a particular computer or device. However, it was found that this was not always necessary as private networks developed and public address space needed to be conserved.
Computers not connected to the Internet, such as factory machines that communicate only with each other via TCP/IP, need not have globally unique IP addresses. Three ranges of IPv4 addresses for private networks were reserved in RFC 1918. These addresses are not routed on the Internet and thus their use need not be coordinated with an IP address registry.
Today, when needed, such private networks typically connect to the Internet through network address translation (NAT).
| Start | End | No. of addresses | |
|---|---|---|---|
| 24-bit block (/8 prefix, 1 × A) | <IP_ADDRESS> | <IP_ADDRESS> | 16777216 |
| 20-bit block (/12 prefix, 16 × B) | <IP_ADDRESS> | <IP_ADDRESS> | 1048576 |
| 16-bit block (/16 prefix, 256 × C) | <IP_ADDRESS> | <IP_ADDRESS> | 65536 |
Any user may use any of the reserved blocks. Typically, a network administrator will divide a block into subnets; for example, many home routers automatically use a default address range of <IP_ADDRESS> through <IP_ADDRESS> (<IP_ADDRESS>/24).
IPv4 address exhaustion
High levels of demand have decreased the supply of unallocated Internet Protocol Version 4 (IPv4) addresses available for assignment to Internet service providers and end user organizations since the 1980s. This development is referred to as IPv4 address exhaustion. IANA’s primary address pool was exhausted on 3 February 2011, when the last five blocks were allocated to the five RIRs. [5] [6] APNIC was the first RIR to exhaust its regional pool on 15 April 2011, except for a small amount of address space reserved for the transition to IPv6, intended to be allocated in a restricted process. [7]
IPv6 addresses
The rapid exhaustion of IPv4 address space prompted the Internet Engineering Task Force (IETF) to explore new technologies to expand the addressing capability in the Internet. The permanent solution was deemed to be a redesign of the Internet Protocol itself. This new generation of the Internet Protocol was eventually named Internet Protocol Version 6 (IPv6) in 1995. The address size was increased from 32 to 128 bits (16 octets), thus providing up to 2 128 (approximately 3.403×10 38 ) addresses. This is deemed sufficient for the foreseeable future.
The intent of the new design was not to provide just a sufficient quantity of addresses, but also redesign routing in the Internet by more efficient aggregation of subnetwork routing prefixes. The resulted in slower growth of routing tables in routers. The smallest possible individual allocation is a subnet for 2 64 hosts, which is the square of the size of the entire IPv4 Internet. At these levels, actual address utilization rates will be small on any IPv6 network segment. The new design also provides the opportunity to separate the addressing infrastructure of a network segment, i.e. the local administration of the segment’s available space, from the addressing prefix used to route traffic to and from external networks. IPv6 has facilities that automatically change the routing prefix of entire networks, should the global connectivity or the routing policy change, without requiring internal redesign or manual renumbering.
The large number of IPv6 addresses allows large blocks to be assigned for specific purposes and, where appropriate, to be aggregated for efficient routing. With a large address space, there is no need to have complex address conservation methods as used in CIDR.
All modern desktop and enterprise server operating systems include native support for the IPv6 protocol, but it is not yet widely deployed in other devices, such as residential networking routers, voice over IP (VoIP) and multimedia equipment, and network peripherals.
Private addresses
Just as IPv4 reserves addresses for private networks, blocks of addresses are set aside in IPv6. In IPv6, these are referred to as unique local addresses (ULA). RFC 4193 reserves the routing prefix fc00::/7 for this block which is divided into two /8 blocks with different implied policies. The addresses include a 40-bit pseudorandom number that minimizes the risk of address collisions if sites merge or packets are misrouted. [8]
Early practices used a different block for this purpose (fec0::), dubbed site-local addresses. However, the definition of what constituted sites remained unclear and the poorly defined addressing policy created ambiguities for routing. This address type was abandoned and must not be used in new systems.
Addresses starting with fe80:, called link-local addresses, are assigned to interfaces for communication on the attached link. The addresses are automatically generated by the operating system for each network interface. This provides instant and automatic communication between all IPv6 host on a link. This feature is required in the lower layers of IPv6 network administration, such as for the Neighbor Discovery Protocol.
Private address prefixes may not be routed on the public Internet.
IP subnetworks
IP networks may be divided into subnetworks in both IPv4 and IPv6. For this purpose, an IP address is logically recognized as consisting of two parts: the network prefix and the host identifier , or interface identifier (IPv6). The subnet mask or the CIDR prefix determines how the IP address is divided into network and host parts.
The term subnet mask is only used within IPv4. Both IP versions however use the CIDR concept and notation. In this, the IP address is followed by a slash and the number (in decimal) of bits used for the network part, also called the routing prefix . For example, an IPv4 address and its subnet mask may be <IP_ADDRESS> and <IP_ADDRESS>, respectively. The CIDR notation for the same IP address and subnet is <IP_ADDRESS>/24, because the first 24 bits of the IP address indicate the network and subnet.
IP address assignment
Internet Protocol addresses are assigned to a host either anew at the time of booting, or permanently by fixed configuration of its hardware or software. Persistent configuration is also known as using a static IP address . In contrast, in situations when the computer’s IP address is assigned newly each time, this is known as using a dynamic IP address .
Methods
Static IP addresses are manually assigned to a computer by an administrator. The exact procedure varies according to platform. This contrasts with dynamic IP addresses, which are assigned either by the computer interface or host software itself, as in Zeroconf, or assigned by a server using Dynamic Host Configuration Protocol (DHCP). Even though IP addresses assigned using DHCP may stay the same for long periods of time, they can generally change. In some cases, a network administrator may implement dynamically assigned static IP addresses. In this case, a DHCP server is used, but it is specifically configured to always assign the same IP address to a particular computer. This allows static IP addresses to be configured centrally, without having to specifically configure each computer on the network in a manual procedure.
In the absence or failure of static or stateful (DHCP) address configurations, an operating system may assign an IP address to a network interface using state-less auto-configuration methods, such as Zeroconf.
Uses of dynamic address assignment
IP addresses are most frequently assigned dynamically on LANs and broadband networks by the Dynamic Host Configuration Protocol (DHCP). They are used because it avoids the administrative burden of assigning specific static addresses to each device on a network. It also allows many devices to share limited address space on a network if only some of them will be online at a particular time. In most current desktop operating systems, dynamic IP configuration is enabled by default so that a user does not need to manually enter any settings to connect to a network with a DHCP server. DHCP is not the only technology used to assign IP addresses dynamically. Dialup and some broadband networks use dynamic address features of the Point-to-Point Protocol.
Sticky dynamic IP address
A sticky dynamic IP address is an informal term used by cable and DSL Internet access subscribers to describe a dynamically assigned IP address which seldom changes. The addresses are usually assigned with DHCP. Since the modems are usually powered on for extended periods of time, the address leases are usually set to long periods and simply renewed. If a modem is turned off and powered up again before the next expiration of the address lease, it will most likely receive the same IP address.
Address autoconfiguration
RFC 3330 defines an address block, <IP_ADDRESS>/16, for the special use in link-local addressing for IPv4 networks. In IPv6, every interface, whether using static or dynamic address assignments, also receives a local-link address automatically in the block fe80::/10.
These addresses are only valid on the link, such as a local network segment or point-to-point connection, that a host is connected to. These addresses are not routable and like private addresses cannot be the source or destination of packets traversing the Internet.
When the link-local IPv4 address block was reserved, no standards existed for mechanisms of address autoconfiguration. Filling the void, Microsoft created an implementation that is called Automatic Private IP Addressing (APIPA). APIPA has been deployed on millions of machines and has, thus, become a de facto standard in the industry. In RFC 3927, the IETF defined a formal standard for this functionality, entitled Dynamic Configuration of IPv4 Link-Local Addresses .
Uses of static addressing
Some infrastructure situations have to use static addressing, such as when finding the Domain Name System (DNS) host that will translate domain names to IP addresses. Static addresses are also convenient, but not absolutely necessary, to locate servers inside an enterprise. An address obtained from a DNS server comes with a time to live, or caching time, after which it should be looked up to confirm that it has not changed. Even static IP addresses do change as a result of network administration (RFC 2072).
Routing
IP addresses are classified into several classes of operational characteristics: unicast, multicast, anycast and broadcast addressing.
Unicast addressing
The most common concept of an IP address is in unicast addressing, available in both IPv4 and IPv6. It normally refers to a single sender or a single receiver, and can be used for both sending and receiving. Usually, a unicast address is associated with a single device or host, but a device or host may have more than one unicast address. Some individual PCs have several distinct unicast addresses, each for its own distinct purpose. Sending the same data to multiple unicast addresses requires the sender to send all the data many times over, once for each recipient.
Broadcast addressing
In IPv4 it is possible to send data to all possible destinations (“all-hosts broadcast”), which permits the sender to send the data only once, and all receivers receive a copy of it. In the IPv4 protocol, the address <IP_ADDRESS> is used for local broadcast. In addition, a directed (limited) broadcast can be made by combining the network prefix with a host suffix composed entirely of binary 1s. For example, the destination address used for a directed broadcast to devices on the <IP_ADDRESS>/24 network is <IP_ADDRESS>. IPv6 does not implement broadcast addressing and replaces it with multicast to the specially-defined all-nodes multicast address.
Multicast addressing
A multicast address is associated with a group of interested receivers. In IPv4, addresses <IP_ADDRESS> through <IP_ADDRESS> (the former Class D addresses) are designated as multicast addresses. IPv6 uses the address block with the prefix ff00::/8 for multicast applications. In either case, the sender sends a single datagram from its unicast address to the multicast group address and the intermediary routers take care of making copies and sending them to all receivers that have joined the corresponding multicast group.
Anycast addressing
Like broadcast and multicast, anycast is a one-to-many routing topology. However, the data stream is not transmitted to all receivers, just the one which the router decides is logically closest in the network. Anycast address is an inherent feature of only IPv6. In IPv4, anycast addressing implementations typically operate using the shortest-path metric of BGP routingand do not take into account congestion or other attributes of the path. Anycast methods are useful for global load balancing and are commonly used in distributed DNS systems.
Public address es
A public IP address , in common parlance, is synonymous with a globally routable unicast IP address .
Both IPv4 and IPv6 define address ranges that are reserved for private networks and link-local addressing. The term public IP address often used excludes these types of addresses.
Modifications to IP addressing
IP blocking and firewalls
Firewalls perform Internet Protocol blocking to protect networks from unauthorized access. They are common on today’s Internet. They control access to networks based on the IP address of a client computer. Whether using a blacklist or a whitelist, the IP address that is blocked is the perceived IP address of the client, meaning that if the client is using a proxy server or network address translation, blocking one IP address may block many individual computers.
IP address translation
Multiple client devices can appear to share IP addresses: either because they are part of a shared hosting web server environment or because an IPv4 network address translator (NAT) or proxy server acts as an intermediary agent on behalf of its customers, in which case the real originating IP addresses might be hidden from the server receiving a request. A common practice is to have a NAT hide a large number of IP addresses in a private network. Only the “outside” interface(s) of the NAT need to have Internet-routable addresses.
Most commonly, the NAT device maps TCP or UDP port numbers on the side of the larger, public network to individual private addresses on the masqueraded network.
In small home networks, NAT functions are usually implemented in a residential gateway device, typically one marketed as a “router”. In this scenario, the computers connected to the router would have private IP addresses and the router would have a public address to communicate on the Internet. This type of router allows several computers to share one public IP address.
Diagnostic tools
Computer operating systems provide various diagnostic tools to examine their network interface and address configuration. Windows provides the command-line interface tools ipconfig and netsh and users of Unix-like systems can use ifconfig, netstat, route, lanstat, fstat, or iproute2 utilities to accomplish the task.
- IP Address. Provided by : Wikipedia. Located at : https://en.wikipedia.org/wiki/IP_address . License : CC BY-SA: Attribution-ShareAlike
- Tiny house move. Authored by : Tammy Strobel. Located at : https://www.flickr.com/photos/rowdykittens/7944094174/ . License : CC BY: Attribution
- route-moutonniere Alger. Authored by : abdelkrim Sa. Located at : https://www.flickr.com/photos/algeria_imaging_pro/16554842499/ . License : CC BY-NC-ND: Attribution-NonCommercial-NoDerivatives


5.1.3C: Role Conflict
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Learning Objectives
- Interpret how role conflict affects an individual’s perception of him/herself and his/her place within society
Role conflict describes a conflict between or among the roles corresponding to two or more statuses fulfilled by one individual. We experience role conflict when we find ourselves pulled in various directions as we try to respond to the many statuses we hold.
The most obvious example of role conflict is work/family conflict, or the conflict one feels when pulled between familial and professional obligations. Take, for example, a mother who is also a doctor. She likely has to work long hours at the hospital and may even be on call several nights a week, taking her away from her children. Many individuals who find themselves in this position describe feeling conflicted and distressed about their situation. In other words, they experience role conflict.
The above example is presented as a personality role conflict: The woman is torn between the part of her personality that values being a mother and the part that identifies as being a doctor. For an example of interpersonal relations causing role conflict, consider an individual who is a school superintendent and a father. He might think that his wife and children expect him to spend most of his evenings with them, but he may also feel that his school board and parents’ groups expect him to spend most of his after-office hours on educational and civic activities. He is therefore unable to satisfy both of these incompatible expectations, and role conflict is the result.
Conflict among the roles begins because of the human’s desire to reach success, and because of the pressure put on an individual by two imposing, apposing and incompatible demands competing against each other. The effects of role conflict, as found through case studies and nationwide surveys, are related to individual personality characteristics and interpersonal relations. Individual personality characteristic conflicts can arise when “aspects of an individual’s personality are in conflict with other aspects of that same individual’s personality. ”
Key Points
- We experience role conflict when we find ourselves pulled in various directions as we try to respond to the many statuses we hold.
- The most common form of role conflict is work/family conflict, in which one needs to prioritize familial or professional obligations.
- The effects of role conflict are related to individual personality characteristics and interpersonal relations.
Key Terms
- work/family conflict : A conflict one faces when one must choose between family needs and work obligations.
- role conflict : A conflict between or among the roles corresponding to two or more statuses in one individual.


3.2: Health History Basic Concepts
During a health history, the nurse collects subjective data from the patient, their caregivers, and/or family members using focused and open-ended questions. Before discussing the components of a health history, let’s review some important concepts related to assessment and communicating effectively with patients.
Subjective Versus Objective Data
Obtaining a patient’s health history is a component of the Assessment phase of the nursing process. Information obtained while performing a health history is called subjective data. Subjective data is information obtained from the patient and/or family members and can provide important cues about functioning and unmet needs requiring assistance. Subjective data is considered a symptom because it is something the patient reports. When documenting subjective data in a progress note, it should be included in quotation marks and start with verbiage such as, “The patient reports…” or “The patient’s wife states…” An example of subjective data is when the patient reports, “I feel dizzy.”
A patient is considered the primary source of subjective data. Secondary sources of data include information from the patient’s chart, family members, or other health care team members. Patients are often accompanied by their care partners. Care partners are family and friends who are involved in helping to care for the patient. For example, parents are care partners for children; spouses are often care partners for each other, and adult children are often care partners for their aging parents. When obtaining a health history, care partners may contribute important information related to the health and needs of the patient. If data is gathered from someone other than the patient, the nurse should document where the information is obtained.
Objective data is information observed through your senses of hearing, sight, smell, and touch while assessing the patient. Objective data is obtained during the physical examination component of the assessment process. Examples of objective data are vital signs, physical examination findings, and laboratory results. An example of objective data is recording a blood pressure reading of 140/86. Subjective data and objective data are often recorded together during an assessment. For example, the symptom the patient reports, “I feel itchy all over,” is documented in association with the sign of an observed raised red rash located on the upper back and chest.
Addressing Barriers and Adapting Communication
It is vital to establish rapport with a patient before asking questions about sensitive topics to obtain accurate data regarding the mental, emotional, and spiritual aspects of a patient’s condition. When interviewing a patient, also consider the patient’s developmental status and level of understanding. Ask one question at a time and allow adequate time for the patient to respond. If the patient does not provide an answer even with additional time, try rephrasing the question in a different way for improved understanding.
If any barriers to communication exist, adapt your communication to that patient’s specific needs.
Cultural Safety
It is important to conduct a health history in a culturally safe manner. Cultural safety refers to the creation of safe spaces for patients to interact with health professionals without judgment or discrimination. Focus on factors related to a person’s cultural background that may influence their health status. It is helpful to use an open-ended question to allow the patient to share what they believe to be important. For example, ask “I am interested in your cultural background as it relates to your health. Can you share with me what is important to know about your cultural background as part of your health care?”
If a patient’s primary language is not English, it is important to obtain a medical translator, as needed, prior to initiating the health history. The patient’s family member or care partner should not interpret for the patient. The patient may not want their care partner to be aware of their health problems or their care partner may not be familiar with correct medical terminology that can result in miscommunication.


15.1: The Elements of Music
15.1 The Elements of Music
Creating contrast usually involves varying one or more of the “Elements of Music,” which are listed and explained below.
| Melody |
We all know what melody is, but don’t confuse melody with ostinato! Melody lives and breathes and has variety. |
| Harmony |
In this section, we will associate harmony with key and mode (major versus minor). |
| Rhythm |
What the most common rhythmic value is in a section. |
| Timbre |
The different tone colors of different instruments, and the different tone colors an individual instrument can produce. |
| Texture |
Described in the previous chapter as existing in the vertical plane (the number of voices) and the horizontal plane (rhythmic activity of the various voices). |
| Articulation |
Staccato versus legato, or short, clipped notes versus long, sustained or connected notes. |
| Dynamics |
Loud versus soft. |
| Register |
High versus middle versus low. |
In the next two sections we will look at how two different pieces exhibit contrast between various sections using the elements of music.


21.1: Unit 1 Overview and Objectives
Overview
As with any journey, it is easier to get started if you prepare appropriately. The assessments in Unit 1 will help you do just that.
Learning Objectives
At the end of unit, you will be able to:
- Determine your compatibility with distance learning
- Identify course related policies and procedures
- Identify health characteristics that affect exercise safety
- Differentiate between physical activity and exercise
- Evaluate your current dimensions of wellness
- List behaviors that contribute to a wellness lifestyle and how your behaviors compare with them
- Describe the stages of behavior change and determine your stage for several wellness-related behaviors
- Identify the strategies in creating an individualized plan to change a wellness-related behavior
CC licensed content, Shared previously
- Lifetime Fitness and Wellness. Provided by : Extended Learning Institute of Northern Virginia Community College. Located at : https://eli.nvcc.edu/ . License : CC BY: Attribution


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2.7: Interest Paid on Bonds and Dividends Paid on Stock
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Last updated
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Save as PDF
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Interest
on debt
is tax deductible
to the corporation
; dividends
on preferred and common
stocks
are not
tax-deductible
, under current law
.
-
Interest must be paid before any dividend payments on preferred and common stocks may be made.
-
Preferred
stock
dividends are usually fixed (like the interest on most bonds).
Even if the corporation becomes more profitable, the preferred stock dividend
can
not
be increased. There is no upside, in this case.
-
Only a
fter interest is paid on the corporation’s debt,
may
preferred dividends then be paid. If the dividend is not paid, it is not consid
ered a “default” as with a loan
or debt; this is because preferred stock represents ownershi
p interests and not
a
liability. Preferred stock is thus thought of as a hybrid debt/equity security
as it has characteristics of both
.
-
Most preferred shares are “cumulative,” which means that before any dividends are paid to common shareholders, all those preferred dividends
that
have
not been paid, and are thus said to have
accumulated
unpaid or
“in arrears
,
”
must first be paid
2
.
-
Common stock is most risky
– first,
because
its dividends are
the
last to be paid and
secondly,
because common shareholders are
the
last to
be paid off in bankruptcy (thus the phrase “residual interest,” used above)
.
-
However,
as common shareholders have rights to “
residual
” profits (i.e., after interest is paid on debt and, second, after preferred stock dividend distributions) that the firm ma
y generate
,
common shareholders
also
have
the most opportunity
to share in positive earnings growth
, i.e., they have the most to gain
.
-
As a result, c
ommon shares usually come with “voting rights,” i.e., the ability annually to vote for company management and on certain key issues. Notably preferred shares rarely carry such rights.


22.3: Relating Angular and Translational Quantities
In this section, we relate each of the rotational variables to the translational variables defined in 1 Dimensional Kinematics and 2 Dimensional Kinematics . This will complete our ability to describe rigid-body rotations.
Angular vs. Linear Variables
In Rotational Variables , we introduced angular variables. If we compare the rotational definitions with the definitions of linear kinematic variables, we find that there is a mapping of the linear variables to the rotational ones. Linear position, velocity, and acceleration have their rotational counterparts, as we can see when we write them side by side:
| Linear | Rotational | |
|---|---|---|
| Position | \(x\) | \(\theta\) |
| Velocity | \(v = \frac{dx}{dt}\) | \(\omega = \frac{d \theta}{dt}\) |
| Acceleration | \(a = \frac{dv}{dt}\) | \(a = \frac{d \omega}{dt}\) |
| Mass 1 | \(m\) | \(I\) |
Let’s compare the linear and rotational variables individually. The linear variable of position has physical units of meters, whereas the angular position variable has dimensionless units of radians, as can be seen from the definition of \(\theta = \frac{s}{r}\), which is the ratio of two lengths. The linear velocity has units of m/s, and its counterpart, the angular velocity, has units of rad/s. In Rotational Variables , we saw in the case of circular motion that the linear tangential speed of a particle at a radius r from the axis of rotation is related to the angular velocity by the relation v t = r\(\omega\). This could also apply to points on a rigid body rotating about a fixed axis. Here, we consider only circular motion. In circular motion, both uniform and nonuniform, there exists a centripetal acceleration ( Motion in Two and Three Dimensions ). The centripetal acceleration vector points inward from the particle executing circular motion toward the axis of rotation. The derivation of the magnitude of the centripetal acceleration is given in Motion in Two and Three Dimensions . From that derivation, the magnitude of the centripetal acceleration was found to be
\[a_{c} = \frac{v_{t}^{2}}{r}, \label{10.14}\]
where r is the radius of the circle.
Thus, in uniform circular motion when the angular velocity is constant and the angular acceleration is zero, we have a linear acceleration—that is, centripetal acceleration—since the tangential speed in Equation \ref{10.14} is a constant. If nonuniform circular motion is present, the rotating system has an angular acceleration, and we have both a linear centripetal acceleration that is changing (because v t is changing) as well as a linear tangential acceleration. These relationships are shown in Figure \(\PageIndex{1}\), where we show the centripetal and tangential accelerations for uniform and nonuniform circular motion.
The centripetal acceleration is due to the change in the direction of tangential velocity, whereas the tangential acceleration is due to any change in the magnitude of the tangential velocity. The tangential and centripetal acceleration vectors \(\vec{a}_{t}\) and \(\vec{a}_{c}\) are always perpendicular to each other, as seen in Figure \(\PageIndex{1}\). To complete this description, we can assign a total linear acceleration vector to a point on a rotating rigid body or a particle executing circular motion at a radius r from a fixed axis. The total linear acceleration vector \(\vec{a}\) is the vector sum of the centripetal and tangential accelerations,
\[\vec{a} = \vec{a}_{c} + \vec{a}_{t} \ldotp \label{10.15}\]
The total linear acceleration vector in the case of nonuniform circular motion points at an angle between the centripetal and tangential acceleration vectors, as shown in Figure \(\PageIndex{2}\). Since \(\vec{a}_{c} \perp \vec{a}_{t}\), the magnitude of the total linear acceleration is
\[|\vec{a}| = \sqrt{a_{c}^{2} + a_{t}^{2}} \ldotp\]
Note that if the angular acceleration is zero, the total linear acceleration is equal to the centripetal acceleration.
Relationships between Rotational and Translational Motion
We can look at two relationships between rotational and translational motion.
- Generally speaking, the linear kinematic equations have their rotational counterparts. Table 10.2 lists the four linear kinematic equations and the corresponding rotational counterpart. The two sets of equations look similar to each other, but describe two different physical situations, that is, rotation and translation.
Table 10.2 - Rotational and Translational Kinematic Equations
| Rotational | Translational |
|---|---|
| \(\theta_{f} = \theta_{0} + \bar{\omega} t\) |
\(x = x_{0} + \bar{v} t\) |
| \(\omega_{f} = \omega_{0} + \alpha t\) | \(v_{f} = v_{0} + at\) |
| \(\theta_{f} = \theta_{0} + \omega_{0} t + \frac{1}{2} at^{2}\) | \(x_{f} = x_{0} + v_{0} t + \frac{1}{2} \omega t^{2}\) |
| \(\omega_{f}^{2} = \omega_{0}^{2} + 2 \alpha (\Delta \theta)\) | \(v_{f}^{2} = v_{0}^{2} + 2a (\Delta x)\) |
- The second correspondence has to do with relating linear and rotational variables in the special case of circular motion. This is shown in Table 10.3, where in the third column, we have listed the connecting equation that relates the linear variable to the rotational variable. The rotational variables of angular velocity and acceleration have subscripts that indicate their definition in circular motion.
Table 10.3 - Rotational and Translational Quantities: Circular Motion
| Rotational | Translational | Relationship ( r = radius |
|---|---|---|
| \(\theta\) | \(s\) | \(\theta = \frac{s}{r}\) |
| \(\omega\) | \(v_{t}\) | \(\omega = \frac{v_{t}}{r}\) |
| \(\alpha\) | \(a_{t}\) | \(\alpha = \frac{a_{t}}{r}\) |
| \(a_{c}\) | \(a_{c} = \frac{v_{t}^{2}}{r}\) |
A centrifuge has a radius of 20 cm and accelerates from a maximum rotation rate of 10,000 rpm to rest in 30 seconds under a constant angular acceleration. It is rotating counterclockwise. What is the magnitude of the total acceleration of a point at the tip of the centrifuge at t = 29.0s? What is the direction of the total acceleration vector?
Strategy
With the information given, we can calculate the angular acceleration, which then will allow us to find the tangential acceleration. We can find the centripetal acceleration at t = 0 by calculating the tangential speed at this time. With the magnitudes of the accelerations, we can calculate the total linear acceleration. From the description of the rotation in the problem, we can sketch the direction of the total acceleration vector.
Solution
The angular acceleration is
\[\alpha = \frac{\omega - \omega_{0}}{t} = \frac{0 - (1.0 \times 10^{4}) \left(\dfrac{2 \pi\; rad}{60.0\; s}\right)}{30.0\; s} = -34.9\; rad/s^{2} \ldotp\]
Therefore, the tangential acceleration is
\[a_{t} = r \alpha = (0.2\; m)(-34.9\; rad/s^{2}) = -7.0\; m/s^{2} \ldotp\]
The angular velocity at t = 29.0 s is
\[\begin{split} \omega & = \omega_{0} + \alpha t = (1.0 \times 10^{4}) \left(\dfrac{2 \pi\; rad}{60.0\; s}\right) + (-39.49\; rad/s^{2})(29.0\; s) \\ & = 1047.2\; rad/s - 1012.71\; rad/s = 35.1\; rad/s \ldotp \end{split}\]
Thus, the tangential speed at t = 29.0 s is
\[v_{t} = r \omega = (0.2\; m)(35.1\; rad/s) = 7.0\; m/s \ldotp\]
We can now calculate the centripetal acceleration at t = 29.0 s:
\[a_{c} = \frac{v^{2}}{r} = \frac{(7.0\; m/s)^{2}}{0.2\; m} = 245.0\; m/s^{2} \ldotp\]
Since the two acceleration vectors are perpendicular to each other, the magnitude of the total linear acceleration is
\[|\vec{a}| = \sqrt{a_{c}^{2} + a_{t}^{2}} = \sqrt{(245.0)^{2} + (-7.0)^{2}} = 245.1\; m/s^{2} \ldotp\]
Since the centrifuge has a negative angular acceleration, it is slowing down. The total acceleration vector is as shown in Figure \(\PageIndex{3}\). The angle with respect to the centripetal acceleration vector is
\[\theta = \tan^{-1} \left(\dfrac{-7.0}{245.0}\right) = -1.6^{o} \ldotp\]
The negative sign means that the total acceleration vector is angled toward the clockwise direction.
Significance
From Figure \(\PageIndex{3}\), we see that the tangential acceleration vector is opposite the direction of rotation. The magnitude of the tangential acceleration is much smaller than the centripetal acceleration, so the total linear acceleration vector will make a very small angle with respect to the centripetal acceleration vector.
A boy jumps on a merry-go-round with a radius of 5 m that is at rest. It starts accelerating at a constant rate up to an angular velocity of 5 rad/s in 20 seconds. What is the distance travelled by the boy?
Check out this PhET simulation to change the parameters of a rotating disk (the initial angle, angular velocity, and angular acceleration), and place bugs at different radial distances from the axis. The simulation then lets you explore how circular motion relates to the bugs’ xy-position, velocity, and acceleration using vectors or graphs.
1 It is a little bit strange to put mass in this table of quantities that are purely related to the motion of the object; however, it is worth pointing out that the moment of inertia (\(I\)) is very much the rotational equivalent of mass (\(m\))